Math 237 Lecture Notes
Math 237 Lecture Notes
Queens University
Mathematics and Engineering and Mathematics and Statistics
Serdar Yuksel
ii
This document is a collection of supplemental lecture notes used for MTHE / MATH 237: Differential Equations for
Engineering Science.
Serdar Yuksel
Contents
1.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
1.2.1
1.2.2
1.2.3
1.2.4
1.2.5
2.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
2.7 Applications . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
3.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
3.2.1
Linear Independence . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 10
3.2.2
iv
Contents
3.2.3
4
Non-Homogeneous Problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 11
4.2.2
4.2.3
4.3.2
Variation of Parameters . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 16
4.3.3
Reduction of Order . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 17
Resonance . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 19
Linear Systems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 29
7.2.2
Non-Linear Systems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 30
Contents
8.2 Transformations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 33
8.2.1
Laplace Transform . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 33
8.2.2
8.2.3
8.2.4
8.2.5
Convolution: . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 35
8.2.6
Step Function . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 35
8.2.7
Impulse response . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 35
10 Series Method . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 39
10.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 39
10.2 Taylor Series Expansions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 39
11 Numerical Methods . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 41
11.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 41
11.2 Numerical Methods . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 41
11.2.1 Eulers Method . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 41
11.2.2 Improved Eulers Method . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 42
11.2.3 Including the Second Order Term in Taylors Expansion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 42
11.2.4 Other Methods . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 42
A
1
Introduction to Differential Equations
1.1 Introduction
Many phenomena in engineering, physics and broad areas of applied mathematics involve entities which change as a
function of one or more variables. The movement of a car along a road, the propagation of sound and waves, the path an
airplane takes, the amount of charge in a capacitor in an electrical circuit, the way a cake taken out from a hot oven and left
in room temperature cools down and many such processes can all be modeled and described as differential equations.
In class we discussed a number of examples such as the swinging pendulum, a falling object with drag and an electrical
circuit involving resistors and capacitors.
Let us make a definition.
Definition 1.1.1 A differential equation is an equation that relates an unknown function and one or more of its derivatives
of with respect to one or more independent variables.
For instance, the equation
dy
= 5x
dx
relates the first derivative of y with respect to x, with x. Here x is the independent variable and y is the unknown function
(dependent variable).
In this course, we will learn how to solve and analyze properties of the solutions to such differential equations. We will
make a precise definition of what solving a differential equation means shortly.
n
d y
Before proceeding further, we make a remark on notation. Recall that dx
n is the n-th derivative of y with respect to x. One
n
d y
(n)
can also use the notation y
to denote dxn . It is further convenient to write y = y (1) and y = y (2) . In physics, the
notation involving dots is also common, such that y denotes the first-order derivative.
Different classifications of differential equations are possible and such classifications make the analysis of the equations
more systematic.
d
1
d2
Q(t) + R 2 Q(t) + Q(t) = E(t),
2
dt
dt
C
which is an equation which arises in electrical circuits. Here, the independent variable is t.
2
f (x, t) = f (x, t)
x2
t
i=0
ay + by (1) + cy = 0
is a linear differential equation, whereas
yy (1) + 5y = 0,
is not linear.
F (x, g(x),
dy dn y
)=0
,
dx dxn
dg(x)
dn g(x)
) = 0,
,...,
dx
dxn
dy dn y
)=0
,
dx dxn
if for all y = g(x) such that G(x, g(x)) = 0, g(x) is an explicit solution to the differential equation on I.
Definition 1.3.3 An n-th parameter family of functions defined on some interval I by the relation
h(x, y, c1 , . . . , cn ) = 0,
is called a general solution of the differential equation if any explicit solution is a member of the family. Each element of
the general solution is a particular solution.
Definition 1.3.4 A particular solution is imposed by supplementary conditions that accompany the differential equations.
If all supplementary conditions relate to a single point, then the condition is called an initial condition. If the conditions
are to be satisfied by two or more points, they are called boundary conditions.
Recall that in class we used the falling object example to see that without a characterization of the initial condition (initial
velocity of the falling object), there exist infinitely many solutions. Hence, the initial condition leads to a particular solution,
whereas the absence of an initial condition leads to a general solution.
Definition 1.3.5 A differential equation together with an initial condition (boundary conditions) is called an initial value
problem (boundary value problem).
The differential equation above tells us what the slope in the graph, that is the ratio of a small change in y with respect to a
small change in x, is. As such, one can draw a graph to obtain the direction field.
Upon a review of preliminary topics on differential equations, we proceed to study first-order ordinary differential equations.
Theorem 1.5.1 Suppose that the real-valued function f (x, y) is defined on a rectangle U = [a, b] [c, d], and suppose
f (x, y) and y
f (x, y) are continuous on U . Suppose further that (x0 , y0 ) is an interior point of U . Then there is an open
subinterval (a1 , b1 ) [a, b], x0 (a1 , b1 ) such that there is exactly one solution to the differential equation
that is defined on (a1 , b1 ) and passes trough the point (x0 , y0 ).
dy
dx
= f (x, y)
The above theorem is important as it tells us that under certain conditions, there can only exist one solution. For example
dy
dy
1/3
might not
dx = sin(y) has a unique solution that passes through a point in the (x, y) plane. On the other hand dx = y
1/3
2/3
have a unique solution in the neighborhood around y = 0, since y y
= (1/3)y
is not continuous when y = 0.
In class, we will have some brief discussion on the existence and uniqueness theorem via the method of successive approximations (which is also known as Piccards method). A curious student is encouraged to think a bit more about the
theorem, on why the conditions are needed.
2
First-Order Ordinary Differential Equations
2.1 Introduction
In this chapter, we consider first-order differential equations. Such a differential equation has the form:
y (1) = f (x, y).
We start the analysis with an important sub-class of first order differential equations. These are exact differential equations.
F (x, y)dx +
F (x, y)dy
x
y
F (x, y)
x
N (x, y) =
F (x, y)
y
and
In this case,
M (x, y)dx + N (x, y)dy = 0,
is called an exact differential equation.
We have the following theorem, to be proven in class:
Theorem 2.2.1 If the functions M (x, y), N (x, y) and their partial derivatives
2
y N (x, y) are continuous on a domain U R , then the differential equation
and
x N (x, y),
M (x, y) =
N (x, y),
y
x
x, y U )
(x, y))M (x, y) + (x, y)( M (x, y)) = ( (x, y))N (x, y) + (x, y)( N (x, y))
y
y
x
x
It is not easy to find such a (x, y) in general. But, if in the above, we further observe that (x, y) is only a function of x,
that is can be written as (x), then we have:
(x)
This can be written as
(
M (x, y) =
(x))N (x, y) + (x) N (x, y)
y
dx
x
d
(x))N (x, y) = (x)( N (x, y)
M (x, y))
dx
x
y
and
If
P (x) =
( y
M (x, y)
x N (x, y))
N (x, y)
(2.1)
is only a function of x, then we can solve the equation (2.1) as a first-order, ordinary differential equation as
(
d
(x)) = (x)P (x),
dx
Rx
P (u)du
(2.2)
K,
1
G(y)f (x) ,
F (x)
g(y)
dx +
dy = 0,
f (x)
G(y)
which is an exact equation. An implicit solution can be obtained by:
Z
Z
g(y)
F (x)
dx +
dy = c,
f (x)
G(y)
One again needs to be cautious with the integrating factor. In particular, if G(y0 ) = 0, then it is possible to lose the solution
y(x) = y0 for all x values. Likewise, if there is an x0 such that f (x0 ) = 0, then, x(y) = x0 for all y is also a solution.
Remark 2.1. If the initial condition is y0 such that G(y0 ) = 0, then y(x) = y0 is a solution.
y
f (x, y) = g( ),
x
for some function g(.).
For differential equations with homogenous coefficients, the transformation y = vx reduces the problem to a separable
equation in x and v. Note that, the essential step is that
dy
d(vx)
dv
=
=v+x .
dx
dx
dx
Alternatively,
M (x, y)dx + N (x, y)dy = 0,
has homogenous coefficients if
M (x, y)
M (x, y)
=
,
N (x, y)
N (x, y)
R, 6= 0
P (x)dx
P (x)dx
(Q(x) P (x)y)dx e
P (x)dx
dy = 0
2.7 Applications
In class, applications in numerous of areas of engineering and applied mathematics will be discussed.
3
Higher-Order Ordinary Linear Differential Equations
3.1 Introduction
In this chapter, we consider differential equations with arbitrary finite orders.
(3.1)
If g(x) = 0 for all x in I, then the differential equation is homogeneous, otherwise, it is non-homogeneous.
For notational convenience, we write the left hand side of (3.1) via:
P(D)(y)(x) = an (x)D(n) (y)(x) + an1 (x)D(n1) (y)(x) + + a1 (x)D(1) (y)(x) + a0 (x)D(0) (y)(x),
which we call as the polynomial operator of order n.
The following is the fundamental theorem of existence and uniqueness applied to such linear differential equations:
Theorem 3.2.1 If a0 (x), a1 (x), . . . , an (x) and g(x) are continuous and real-valued functions of x on an interval I R
and an (x) 6= 0, x I, then the differential equation has a unique solution y(x) which satisfies the initial condition
y(x0 ) = y0 , y (1) (x0 ) = y1 , . . . , y (n1) (x0 ) = yn1 ,
where x0 I, y0 , y1 , . . . , yn1 R.
The following is an immediate consequence of the fundamental theorem:
Theorem 3.2.2 Consider a linear differential equation defined on I:
P(D)(y(x)) = 0,
10
x I and a0 (x), a1 (x), . . . , an (x) and g(x) are continuous and real-valued functions of x with
y(x0 ) = 0, y (1) (x0 ) = 0, . . . , y (n1) (x0 ) = 0,
I.
Then,
x I
y(x) = 0,
Exercise 3.2.1 Let f1 , f2 , . . . , fn be solutions to
P(D)(y(x)) = 0.
Then,
c1 f1 (x) + c2 f2 (x) + + cn fn (x) = 0,
is also a solution, for all c1 , c2 . . . , cn R.
3.2.1 Linear Independence
Linear independence/dependence is an important concept which arises in differential equations, as well as in linear algebra.
Definition 3.2.2 The set of functions {f1 (x), f2 (x), f3 (x), . . . , fn (x)} are linearly independent on an interval I if there
exist constants, c1 , c2 , . . . , cn ; not all zero, such that
c1 f1 (x) + c2 f2 (x) + + cn fn (x) = 0,
x I.
W (f1 , f2 , . . . , fn ) = det
f1
(1)
f1
...
...
(n1)
f1
f2
. . . fn
(1)
(1)
. . . , fn
f2
... ... ...
11
Theorem 3.2.4 Let f1 , f2 , . . . , fn be n 1st order differentiable functions on an interval I. If these functions are linearly
dependent, then the Wronskian is identically zero on I.
The proof of this is left as an exercise. The student is encouraged to use the definition of linear independence to show that
the Wronskian is zero.
The main use of Wronskian is given by the following, which allows us to present an even stronger result. It says that if the
Wronskian is zero at any given point, then the set of solutions are linearly dependent.
Theorem 3.2.5 Let y 1 , y 2 , . . . , y n be n solutions to the linear equation.
P(D)(y(x)) = 0.
They are linearly dependent if and only if the Wronskian is zero on I. The Wronskian of n solutions to a linear equation is
either identically zero, or is never 0.
The proof of this for a second-order system will be presented as an exercise.
4
Higher-Order Ordinary Linear Differential Equations
4.1 Introduction
In this chapter, we will develop methods for solving higher-order linear differential equations.
i = 1, 2, . . . , n
14
15
Hence, any solution to a particular equation can be obtained by obtaining one particular solution to the non-homogenous
equation, and then adding the solutions to the homogenous equation, and finally obtaining the coefficients for the
homogenous-part of the solution.
Theorem 4.3.1 Let yp be a particular solution to a differential equation
P(D)(y(x)) = g(x),
(4.1)
and let yc = c1 y1 (x) + c2 y2 (x) + + cn yn (x), be the general solution corresponding to the homogeneous equation
P(D)(y(x)) = 0
Then, the general solution of (4.2) is given by
y(x) = yp (x) + yc (x)
Here yc (x) is called the complementary solution.
We now present a more general result. This is called the principle of superposition.
Theorem 4.3.2 Let ypi (x) be respectively particular solutions of
P(D)(y(x)) = gi (x),
(4.2)
for i = 1, 2 . . . , m. Then,
a1 yp1 (x) + a2 yp2 (x) + + am ypm (x)
is a particular solution to the DE
P(D)(y(x)) = a1 g1 (x) + a2 g2 (x) + + am gm (x)
In the following, we discuss how to obtain particular solutions. We first use the method of undetermined coefficients and
then the method of variation of parameters.
The method of undetermined coefficients has limited use, but when it works, is very effective. The method of variation of
parameters is more general, but requires more steps in obtaining solutions.
16
for some constants a, b, p0 , p1 , . . . , pm , q0 , q1 , . . . , qm . If a + ib is not a solution to L(m) = 0, then there exists a particular
solution of the form
yp (x) = eax cos(bx)(Ao + A1 x + A2 x2 + + Am xm ) + eax sin(bx)(Bo + B1 x + B2 x2 + + Bm xm ),
where A0 , A1 , . . . , Am , B0 , B1 , . . . , Bm can be obtained by substitution.
If a + ib is a root of the polynomial with multiplicity k, then the assumed particular solution should be modified by
multiplying with xk .
4.3.2 Variation of Parameters
A particularly effective method is the method of variation of parameters.
Consider,
an (x)y (n) + an1 y (n1) + + a1 (x)y (1) + a0 y = g(x),
where all of the functions are continuous on some interval, where an (x) 6= 0.
Suppose
y(x) = c1 y1 (x) + c2 y2 (x) + + cn yn (x)
is the general solution to the corresponding homogeneous equation. This method, replaces ci with some function ui (x).
That is, we look for u1 (x), u2 (x), . . . , un (x) such that
yp (x) = u1 (x)y1 (x) + u2 (x)y2 (x) + + un (x)yn (x)
This method, due to Lagrange, allows a very large degree of freedom on how to pick the functions u1 , , un . One
restriction is the fact that the assumed solution must satisfy the differential equation. It turns out that, through an intelligent
construction of n 1 other constraints, we can always find such functions: As there are n unknowns, we need n equations,
for which can have freedom on how to choose. We will find these equations as follows:
T
Let U (x) = u1 (x) u2 (x) . . . un (x) and Y (x) = y1 (x) y2 (x) . . . yn (x) . Let us write in vector inner
product form:
17
an Y (n1) U = g(x)
Hence, we obtain:
y1
y1
..
.
(n1)
y1
y2
y2
..
.
(n1)
y2
...
...
yn
yn
..
.
...
(n1)
. . . yn
u1
u2
.. =
.
un
0
0
..
.
g/an
You may recognize that the term above is the matrix M , whose determinant is the Wronskian. We already know that this
matrix is invertible, since the functions y1 , y2 , . . . , yn are linearly independent solutions of the corresponding homogeneous
differential equation.
u1
u2
It follows that, we can solve for . by
..
un
0
0
..
.
(M (y1 , y2 , . . . , yn )(x))1
g(x)
an (x)
u1 (x)
Z
u2 (x)
.. = (M (y1 , y2 , . . . , yn )(x))1
.
un (x)
0
0
..
.
g(x)
an (x)
dx + C,
for some constant C (this constant will also serve us in the homogeneous solution).
Hence, we can find the functions u1 (x), u2 (x), . . . , un (x), and using these, we can find
yp (x) = u1 (x)y1 (x) + u2 (x)y2 (x) + + un (x)yn (x)
As we observed, if n linearly independent solutions to the corresponding homogeneous equation are known, then a particular solution can be obtained for the non-homogenous equation. One question remains however, on how to obtain the
solutions to the homogeneous equation. We now discuss a useful method for obtaining such solutions.
(2y1 + py1 )
v =0
y1
18
This is a first-order equation in v . Thus, v can be solved, leading to a solution for v, and ultimately solving for y(x) =
v(x)y1 (x).
As such, if one solution to the homogeneous equation is known, another independent solution can be obtained.
The same discussion applies for an nth order linear differential equation. If one solution is known, then by writing y(x) =
v(x)y1 (x) and substituting this into the equation (as we did above), a differential equation of order n 1 for v (x) can be
obtained. The n 1 linearly independent solutions for v can all be used to recover n 1 linearly independent solutions to
the original differential equation. Hence, by knowing only one solution, n linearly independent solutions can be obtained.
dqt(t)
dt ,
This is a second order differential equation with constant coefficients. The equation is non-homogeneous. As such, we first
need to obtain the solution to the corresponding homogeneous equation.
This writes as:
m
R
m+
= 0,
L
LC
R
m1 =
+
2L
with solutions:
or
m1 = +
with =
R
2L
and w0 =
1
LC
1
R 2
)
2L
LC
q
(2 w02
q
(2 w02
2 > w02 . In this case, we have two distinct real roots: The general solution is given as qc (t) = c1 em1 t + c2 em2 t
2 = w02 . In this case, we have two equal real roots: The general solution is given as qc (t) = c1 et + c2 tet
19
p
2 = w02 . p
In this case, we have two complex valued roots: The general solution is given as qc (t) = c1 et cos( (2 w02 t)+
c2 et sin( (2 w02 t)
We now can solve the non-homogeneous problem by obtaining a particular solution and adding the complimentary solution
above to the particular solution.
Let E(t) = E cos(wt). In this case, we can use the method of undetermined coefficients to obtain a particular solution.
Suppose w 6= w0 .
Using the method of undetermined coefficients, we find that a particular solution is given by
qp (t) = A cos(wt) + B sin(wt),
where
A=
E
1
2
L w02 w + (2w)
2
2
w0 w
and
B=
2Aw
w02 w2
4.4.1 Resonance
Now, let us consider the case when = 0 and w = w0 , that is the frequency of the input matches the frequency of the
homogeneous equation solution.
In this case, when we use the method of undetermined coefficients, we need to multiply our candidate solution by t to
obtain:
qp (t) = At cos(w0 t) + Bt sin(w0 t),
Substitution yields,
A = 0, B =
E
2w0 L
E
t sin(w0 t)
2w0 L
As can be observed, the magnitude of qp (t) grows over time. This leads to breakdown in many physical systems. However,
if this phenomenon can be controlled (say by having a small non-zero value), resonance can be used for important
applications.
5
Systems of First-Order Linear Differential Equations
5.1 Introduction
In this chapter, we investigate solutions to systems of differential equations.
Definition 5.1.1 A set of differential equations which involve more than one unknown functions and their derivatives with
respect to a single independent variable is called a system of differential equations.
For example
(1)
y1 = f1 (x, y1 , y2 , . . . , yn )
(1)
y2 = f2 (x, y1 , y2 , . . . , yn )
and up to:
yn(1) = fn (x, y1 , y2 , . . . , yn )
is a system of differential equations.
Clearly, the first order equation that we discussed earlier in the semester of the form
y (1) = f (x, y),
is a special case. If each of the functions {f1 , f2 , . . . , fn } is a linear function of {x1 , x2 , . . . , xn }, then the system of
equations is said to be linear. In this case, we have:
(1)
x1 = f1 (t, x1 , x2 , . . . , xn )
(1)
x2 = f2 (t, x1 , x2 , . . . , xn )
and up to:
22
x(1)
n = fn (t, x1 , x2 , . . . , xn )
fi
Theorem 5.1.1 Let { x
}, for all i {1, 2, . . . , n}, j {1, 2, . . . , n}, exist and be continuous in an n + 1-dimensional
j
domain containing the point (t0 , x01 , x02 , . . . , x0n ). Then, there exists an interval [t0 h, t0 + h] with h > 0, such that for
all t [t0 h, t0 + h], there exists a unique solution
x1 (t) = 1 (t),
For a vector x(t) Rn , x (t) denotes the derivative of the vector x(t), and it exists when all the components of x(t) are
differentiable. The derivative of x(t) with respect to t is a vector consisting of the individual derivatives of the components
of x:
x1 (t)
x2 (t)
x (t) = .
..
xn (t)
We note that, the integral of a vector is also defined in a similar pattern, that is, the following holds:
R
R x1 (t)dt
Z
x2 (t)dt
x(t)dt =
..
R .
xn (t)dt
5.1.1 Higher-Order Linear Differential Equations can be reduced to First-Order Systems
One important observation is that any higher-order linear differential equation
y (n) + an1 (x)y (n1) + an2 (x)y (n2) + + a1 (x)y (1) + a0 (x)y = 0,
can be reduced to a first-order system of differential equation by defining:
x1 = y,
x2 = y ,
x3 = y , . . . , xn = y (n1)
It follows that,
x1 = y
x2 = y = x1
x3 = y = x2
until
xn = xn1
As such, we obtain:
0
1
0
x1
x2 0
0
1
x3 0
0
0
=
.. ..
..
..
. .
.
.
a0 (x) a1 (x) a2 (x)
xn
x1
x2
x3
..
.
...
xn
. . . an1 (x)
...
...
...
0
0
0
..
.
23
Hence, if we write:
x1
x2
x = x3
..
.
xn
0
0
0
..
.
1
0
0
..
.
0
1
0
..
.
...
...
...
0
0
0
..
.
A(x) =
...
a0 (x) a1 (x) a2 (x) . . . an1 (x)
We obtain a first-order differential equation:
x (t) = A(x)x(t)
As such, first-order systems of equations are very general.
Exercise 5.1.1 Express 5y (3) + y (2) + y (1) + y = 0 as a system of differential equations, by defining x1 = y, x2 = y and
x3 = y .
Definition 5.2.1 A linearly independent set of solutions to the homogeneous equation is called a fundamental set of solutions.
Theorem 5.2.2 Let {x1 , x2 , . . . , xn } be a fundamental set of solutions to
x (t) = A(t)x(t)
in an interval < t < . Then, the general solution x(t) = (t) can be expressed as
c1 x1 (t) + c2 x2 (t) + . . . cn xn (t),
and for every given initial set of conditions, there is a unique set of coefficients {c1 , c2 , . . . , cn }.
Hence, the main issue is to obtain a fundamental set of solutions. Once we can obtain this, we could obtain the complete
solution to a given, homogeneous differential equation.
24
We will observe that, we will be able to use the method of variation of parameters for systems of equations as well. The
fundamental set of solutions will be useful for this discussion as well.
The next topic will focus on system of differential equations which are linear, and constant-coefficient.
6
Systems of First-Order Constant Coefficient Differential Equations
6.1 Introduction
In this chapter, we consider systems of differential equations with constant coefficients.
t2
tn
+ . . . An + . . .
2
n!
26
10
A=I=
01
In this case,
eAt = I + It + I 2
t2
tn
+ + In + . . .
2!
n!
At
et 0
=
0 et
we obtain
1 0 0
A = 0 2 0 ,
0 0 3
eAt
e 1 0 0
= 0 e2 t 0 ,
0 0 e3 t
Hence, it is very easy to compute the exponential when the matrix has a nice form.
What if A has a Jordan form? We now discuss this case.
First, we use a result that if AB = BA, that is if A and B commute, then
e(A+B) = eA eB
You will prove this in your assignment. In this case, we can write a matrix
1 1 0
A = 0 1 1
0 0 1
as B + C, where
1 0 0
B = 0 1 0
0 0 1
010
C = 0 0 1
000
We note that BC = CB, for B is the identity matrix multiplied by a scalar number. Hence,
eAt = eBt eCt .
All we need to compute is eCt , as we have already discussed how to compute eBt .
27
t2
t3
+ C3 + . . . ,
2!
3!
1 t t2 /2
t
= I + Ct + C 2 = 0 1 t
2!
00 1
2
2
e 1 0 0
1 t t2 /2
e1 t te1 t t2 e1 t
= 0 e1 t 0 0 1 t = 0 e1 t te1 t
00 1
0 0 e1 t
0
0
e1 t
Now that we know how to compute the exponential of a Jordan form, we can proceed to study a general matrix. Let
A = P BP 1 ,
where B is in a Jordan form. Then,
A2 = (P BP 1 )2 = P BP 1 P BP 1 = P B 2 P
A3 = (P BP 1 )3 = P BP 1 P BP 1 P BP 1 = P B 3 P
and hence,
Ak = (P BP 1 )3 = P BP 1 P BP 1 (P BP 1 )k1 = P B k P
Finally,
eA = P (eB )P 1
and
eAt = P (eBt )P 1
Hence, once we obtain a diagonal matrix or a Jordan form matrix B, we can compute the exponential eAt very efficiently.
eA(t ) g( )d
x(t) = eAt x0 +
You could verify this result by substitution. The uniqueness theorem reveals that this has to be the solution.
This equation above is a fundamentally important one for mechanical and control systems. Suppose a spacecraft needs to
move from the Earth to the moon. The path equation is a more complicated version of the equation above, and g(t) is the
control term.
28
7
Stability and Lyapunovs Method
7.1 Introduction
In many engineering applications, one wants to make sure things behave nicely in the long-run; without worrying too much
about the particular path the system takes (so long as these paths are acceptable).
Stability is the characterization ensuring that things behave well in the long-run. We will make this statement precise in the
following.
7.2 Stability
Before proceeding further, let us recall the l2 norm: For a vector x Rn ,
v
u n
uX
x2i ,
||x||2 = t
i=1
30
x = Ax,
the solution is locally and globally asymptotically stable if and only if
max{Re{i }} < 0,
i
where Re{.} denotes the real part of a complex number, and i denotes the eigenvalues of A.
We could also further strengthen the theorem.
Theorem 7.2.2 For a linear differential equation
x = Ax,
the system is locally stable if and only if
max{Re{i }} 0,
i
If Re{i } = 0, for some i , the algebraic multiplicity of this eigenvalue should be the same as the geometric multiplicity.
Here, Re{.} denotes the real part of a complex number, and i denotes the eigenvalues of A.
Exercise 7.2.1 Prove the theorem.
31
b) For a given differential equation x (t) = f (x(t)) with f (0) = 0, and continuous f , if we can find a Lyapunov function
V (x) such that
d
V (x(t)) < 0,
dt
for x(t) = x \ {0} where is a closed and bounded set containing the origin, the system is locally asymptotically
stable.
Proof. a) Let > 0 be given. We will show the existence of > 0 such that for all ||x(0)||2 , ||x(t)||2 for all
t 0. Define m = minx:||x||2= V (x) (such an m exists by the continuity of V ). Let 0 < < so that for all ||x||2 it
follows that V (x) < m. Such a exists by continuity of V . Now, for any ||x(0)||2 , V (x(0)) < m and furthermore by
d
V (x(t)) 0, V (x(t)) < m. This then implies that ||x(t)||2 < since otherwise, there would have to be
the condition dt
some time t1 so that ||x(t1 )||2 = which would result in V (x(t1 )) m, leading to a contradiction.
b) Let be so that with = r, satisfy the condition in part a so that with ||x(0)||2 , ||x(t)||2 r for all t so that
d
V (x(t)) < 0, V (x(t)) is a monotonically decreasing family of non-negative numbers and
x(t) . Since we have that dt
it therefore has a limit. Call this limit c. We will show that c = 0. Suppose c 6= 0. Let be such that for
2 ,
all ||x||
V (x) c. It then follows that for all t, x(t) {x : ||x||2 r}. Define a = max{x:||x||2r} x V (x) f (x).
The number a exists since the functions considered are continuous, and it is either zero or a negative number, but it
cannot be zero by the assumptions stated in the theorem. Therefore, a < 0. This implies then that V (x(t)) = V (x(0)) +
Rt d
s=0 ds V (x(s))ds V (x(0))+at, which implies that V (x(t)) will be negative after a finite t. This cant be true, therefore
c cannot be non-zero.
The theorem above can be strengthened to global asymptotic stability if part b of the Theorem above is also satisfied by
some Lyapunov function with the following property:
1. The Lyapunov function V (x) is radially unbounded, that is lim||x||2 V (x) = and
d
dx V
Exercise 7.3.1 Show that x = x3 is locally asymptotically stable, by picking V (x) = x2 as a Lyapunov function. Is this
solution globally asymptotically stable?
Exercise 7.3.2 Consider x + x + x = 0. Is this system asymptotically stable?
Hint: Convert this equation into a system of first-order differential equations, via x1 = x and x2 = x1 , x2 = x1 x2 .
Then apply V (x1 , x2 ) = x21 + x1 x2 + x22 as a candidate Lyapunov function.
Linearization
Consider x = f (x), where f is a function which has continuous partial derivatives so that it locally admits a first-order
Taylors expansion representation in the following form:
f (x) = f (x0 ) + f (x0 )(x x0 ) + h.o.d.,
where h.o.d. stands for higher order terms. Now consider x0 = 0 and further with the condition that f (0) = 0 and f is
continuously differentiable so that we can write
f (x) = f (0) + f (0)(x) + b(x)x,
where limx0 b(x) = 0. One then studies the properties of the linearized system by only considering the properties of
f (x) at x = 0. Then, if f (0) is a matrix with all of its eigenvalues having negative real parts, then the system is locally
stable.
32
However, it should be noted that, this method only works locally, and hence, the real part of the eigenvalues being less than
zero only implies local stability due to the Taylor series approximation.
You will encounter many applications where stability is a necessary prerequisite for acceptable performance. In particular,
in control systems one essential goal is to adjust the system dynamics such that the system is either locally or globally
asymptpotically stable.
For the scope of this course, all you need to know are the notions of stability and application of Lyapunov theory to
relatively simple systems. The above is provided to provide a more complete picture for the curious.
If time permits, we will discuss these issues further at the end of the semester, but this is highly unlikely. Mathematics and
Engineering students will revisit some of these discussions in different contexts, but with essentially same principles, in
MATH 332, MATH 335 and MATH 430.
8
Laplace Transform Method
8.1 Introduction
In this chapter, we provide an alternative approach to solve constant coefficient differential equations. This method is
known as the Laplace Transform Method.
8.2 Transformations
A transformation T is a map from a signal set to another one, such that T is onto and one-to-one (hence, bijective).
Let us define a set of signals on the time-domain. Consider the following mapping:
Z
K(s, t)x(t)dt, s R.
g(s) =
=
34
for some finite M < and > 0 and f (t) is a piece-wise continuous function in every interval [0, ], 0, R. We
say a function is piece-wise continuous if every finite interval can be divided into a finite number of subintervals in which
the function is continuous, and the function has a finite left-limit and a right-limit in the break-points.
dn
F (s)
dsn
= B(s).
(8.1)
As such, we could obtain the Laplace transform of y(t). We now will have to find the function which has its Laplace
transform as Y (s). Thus, we need to find an inverse Laplace transform.
8.2 Transformations
35
8.2.5 Convolution:
Let f and g be two functions that are piecewise continuous on every finite closed interval, and be less that M et , t 0.
and f (t) = g(t) = 0 for t 0. The function
Z
f ( )g(t )d
h(t) = (f g)(t) =
0
ua (t) = 1(a0) ,
where 1(.) is the indicator function. We can show that for s > 0
L(ua (t)) =
1 as
e
s
if 0 t ,
else.
Now, consider lim0 F (t). This limit is such that, for every t 6= 0, the limit is zero valued. When t = 0, however, the
limit is . of We denote this limit by (t).
This limit is known as the impulse. The impulse is not a regular function. It is not integrable. The proper way to study such
a function is via distribution theory. This will be further covered in Math 334 and Math 335.
For the purpose of this course however, we will always use the impulse function under the integral sign.
In class, we showed that
lim
est F (t)dt = 1,
36
with y(0) = y (1) (0) = = y (n1) (0) = 0. We had observed that the Laplace transform of the solution Y (s) =
G(s)
an sn +an1 sn1 ++a0
1
.
an sn + an1 sn1 + + a0
We denote y(t) when g(t) = (t), by h(t) and call it the impulse response of a system governed by a linear differential
equation.
Exercise 8.2.1 Consider the linear equation above. If g(t) = p(t) for some arbitrary function p(t) defined on t 0, then
the resulting solution is equal to
y(t) = p(t) h(t)
9
On Partial Differential Equations and Separation of Variables
Recall that a differential equation that entails more than one independent variable is called a partial differential equation.
An important example is the heat equation
u
2u
= k 2,
t
x
(9.1)
where u is the dependent function and x, t are the independent variables. Here, k is a constant.
This equation models the variation of temperature u with position x and time t in a heated rod where it is assumed that at
a given cross-section at length x from one end of the rod the temperature is a constant.
Suppose that the rod has a finite length L, extending from x = 0 to x = L. Its temperature function u(x, t) solves the heat
equation. Suppose that the initial temperature is given by
u(x, 0) = f (x)
(9.2)
One can also restrict fixed temperatures at the two ends of the rod so that
u(0, t) = u(L, t) = 0.
(9.3)
A physical setting would be the case where the rod is connected to large objects with constant temperatures.
These equations define a boundary value problem, where (9.1) is defined for 0 < x < L, t > 0.
This example will be studied in this chapter; the ideas used here are applicable to a large class of such differential equations.
First, as in the ordinary differential equations, due to linearity, the principle of superposition holds (see Chapter 4): That
is, if u1 and u2 solve (9.1), so does c1 u1 + c2 u2 for any scalars c1 , c2 . Furthermore, the homogenous conditions (9.3)
are also satisfied by the linear combinations. On the other hand, the condition (9.2) is a non-homogenous condition. If
u(x) = c1 u1 (x, 0) + c2 u2 (x, 0) = f (x) and u1 and u2 satisfy (9.1) and(9.3), then u is a solution to the heat equation given
the boundary conditions.
The solution technique is then to consider first the homogenous problem, and then arrange the terms so as to arrive at the
non-homogenous problem. A common first attempt is to look for a solution that has the structure that:
u(x, t) = X(x)T (t),
where X only depends on X and T only depends on T . Such a method is called the method of separation of variables.
Assuming such a structure holds, (9.1) reduces to:
T (1)
X (2)
=
,
X
kT
38
Now, suppose that both of these terms are equal to a constant . Of course, there may be multiple such values, and we
will need to account for all such values. First, consider the term:
X (2)
+=0
X
We know from our earlier analysis that the general solution for such a problem will be given by either an exponential or a
harmonic family, depending on whether 0 or < 0. The conditions u(0, t) = u(L, t) = 0 actually rule out the case
with 0 for a non-zero solution, and it turns out that the general solution will write as:
where { nL2 , n Z+ }. Now with this set of valus, we can solve the second equation:
n2 2
T (1)
,
=
kT
L2
which leads to the solution:
Tn (t) = Kn e
n2 2 k
t
L2
n2 2 k
t
L2
sin(
n
x)
L
Using the principle of superposition, we can take a linear combination so that any such linear combination
u(x, t) =
X
n
Kn e
n2 2 k
t
L2
sin(
n
x)
L
is a solution to the homogenous problem. Now, if we also satisfy the initial condition:
X
n
u(x, 0) = f (x, 0) =
Kn sin( x),
L
n
given by the non-homogenous boundary condition, we can find the {Kn , n Z+ } coefficients and obtain a solution to the
heat equation.
Next year, in MTHE 334 and MTHE 335; we will see that any continuous function defined one a bounded interval can
be expressed in terms of sines and cosines, through an expansion known as the Fourier Transform. This will allow for the
combination of the Kn terms above, leading to a solution to the heat equation.
Partial differential equations are significantly more challenging in general than ordinary differential equations. Our discussion in this chapter is introductory. However, the method of separation of variables is an effective solution technique
applicable to a large class of problems.
10
Series Method
10.1 Introduction
In this chapter, we discuss the series method to solve differential equations. This is a powerful technique essentially building
on the method of undetermined coefficients.
11
Numerical Methods
11.1 Introduction
In this chapter, we discuss numerical methods to compute solutions to differential equations. Such methods are useful when
a differential equation does not belong to the classifications we have investigated so far in the course.
d
(t t0 )2
d2
x(t)|t=t0 (t t0 ) + 2 x(t)|t=t0
+ h.o.t.,
dt
dt
2!
and
(tk ) = (tk1 ) + f ((tk1 ), tk1 )(tk tk1 )
This method is known as the Eulers method.
42
11 Numerical Methods
1
(tk + h) = (tk ) + x (tk )h + x (tk )h2
2
That is, in the above, we truncate the Taylor expansion at the second term, instead of the first one.
Here one has to compute the second derivative. This can be done as follows:
x =
d dx
d
x
x dx
x
x
d2 x
=
=
x
=
+
=
+
x
dt2
dt dt
dt
t
y dt
t
x
A
Brief Review of Complex Numbers
A.1 Introduction
The square root of -1 does not exist in the space of real numbers. We extend the space of numbers to the so-called complex
ones, so that all polynomials have roots.
We could define the space of complex numbers C as the space of pairs (a, b), where a, b R, on which the following
operations are defined:
(a, b) + (c, d) = (a + b, c + d)
(a, b).(c, d) = (ac bd, ad + bc)
We could also use the notation:
(a, b) = a + bi.
In the above representation, the notation i stands for imaginary. As such b is the imaginary component of the complex
number. With this definition, the square root of -1 is (0, 1) and this is denoted by the symbol i (some texts use the term j,
in particular the literature in physics and electrical engineering heavily use j, instead of i). As such, complex number has
the form
a + bi,
where a denotes the real term, and b is the imaginary term in the complex number. Two complex numbers are equal if and
only if both their real parts and imaginary parts are equal. We denote the space of complex numbers as C.
One could think of a complex number as a vector in a two dimensional space; with the x-axis denoting the real component,
and the imaginary term living on the y-axis. Observe that with b = 0, we recover the real numbers R.
We define the conjugate of a complex number a + bi, is a bi. The absolute value of a complex number is the Euclidean
norm of the vector which represents the complex number, that is
p
|a + bi| = a2 + b2 ,
Complex numbers adopt the algebraic operations of addition, subtraction, multiplication, and division operations that
generalize those of real numbers. These operations are conveniently carried out when the complex number is represented
in terms of exponentials.
44
ex = 1 + x +
xk
x2
+ +
+ ...
2!
k!
z2
zk
+ +
+ ...
2!
k!
In particular, if we assume the above holds for complex numbers as well, it follows that:
(i)2
(i)k
+ +
+ ...
2!
k!
2
3
4
(i)k
= 1 + i
i +
+
+ ...
2!
3!
4!
k!
ei = 1 + i +
(A.1)
Let us recall the Taylor series for cos() and sin() around 0:
cos() = 1
sin() =
2
4
6
2n
+
+ + (1)n
+ ...
2
4!
6!
(2n)!
5
7
2n+1
3
+
+ + (1)n
+ ...
3
5!
7!
(2n + 1)!
The above hold for all R. This can be verified by the fact that series converges and the error in the approximation
converges to 0.
Thus, ei satisfies:
ei = cos() + i sin()
The above is known as Eulers formula and is one of the fundamental relationships in applied mathematics. Also, note that
via the equations above, it follows that
ei + ei
,
cos() =
2
and
ei ei
sin() =
2i
Now, we could represent complex numbers in terms of exponentials. For example
ei/2 = i
ei = 1
ei3/2 = i
ei2 = 1
p
a2 + b 2
tan() = (b/a),
since a = a2 + b2 cos() and b = a2 + b2 sin() Such an exponential representation of a complex number also lets us
take roots of complex numbers: For example let us compute:
45
(1)1/4
This might arise for example, when one tries to solve the homogenous differential equation
y (4) + y = 0
(A.2)
1
1
1
1
y(x) = c1 cos( x) + c2 sin( x) + c3 cos(3 x) + c4 sin(3 x).
4
4
4
4
B
Similarity Transformation and the Jordan Canonical Form
B.2 Diagonalization
If A is a real symmetric matric, and if we choose P to be a matrix whose columns are the orthonormal eigenvectors of A,
then the matrix B = P 1 AP is diagonal, which is also real, by an argument we proved in class. Hence, real symmetric
matrices can be diagonalized by a similarity transformation.
However, if the matrix A is not symmetric, then diagonalization might not be possible.
For such systems, a structure, which is close to a diagonal form, known as the Jordan canonical form, can be obtained:
B1 0 . . . 0
0 ...
0 B2 . . . 0
0 ...
.. .. .. ..
..
..
. . . .
.
.
B=
0 0 . . . Bk 0 . . .
0 0 . . . 0 Bk+1 . . .
.. .. .. ..
..
..
. . . .
.
.
0
0 0
0
0
..
.
..
.
. . . Bm
48
i
0
Bi = .
..
1
i
..
.
0 ...
1 ...
.. ..
. .
0
0
.. .
.
0 0 0 . . . i
One could also have a Jordan canonical matrix with only real-valued entries (but we will not consider this).
In class we defined the notions of algebraic multiplicity and geometric multiplicity of an eigenvalue. Recall that the algebraic multiplicity of an eigenvalue i is equal to the number of times i appears as a root to det(A I) = 0.
The geometric multiplicity is the size of N (A I) (the null space of A I). This is equal to the number of linearly
independent eigenvectors that can be obtained as a result of the equation: Ax = i x.
If the matrix A has at least one eigenvalue with its algebraic multiplicity greater than its geometric multiplicity, then
diagonalization is not possible (Hence, it follows that, having a repeated eigenvalue does not automatically mean that
diagonalization is not possible).
A convenient way to generate a Jordan canonical form is through the use of eigenvectors and generalized eigenvectors:
That is, the columns of the matrix P may consist of the eigenvectors and the generalized eigenvectors.
Clearly, any n n square matrix with distinct eigenvalues (with n different eigenvalues) can be transformed into a diagonal
matrix, as there will be n Jordan blocks which form the diagonal matrix.