Iit - Foundation - Set - V
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Hargovind NainaniIit - Foundation - Set - V
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IIT FOUNDATION
SET - 5
INDEX
PHYSICS
1. RIGID BODY DYNAMICS ..................................................................................
2
2. SIMPLE HARMONIC MOTION .......................................................................
24
3. SOUND ................................................................................................................
42
4. MAGNETISM .....................................................................................................
70
5. ALTERNATING CURRENTS .............................................................................
94
6. ELECTRO MAGNETISM .................................................................................. 114
7. THERMO ELECTRICITY ................................................................................. 150
8. SEMICONDUCTOR ........................................................................................... 174
9. RAY OPTICS ....................................................................................................... 184
10. PHYSICAL OPTICS ........................................................................................... 210
CHEMISTRY
1. ATOMIC STRUCTURE ...................................................................................... 233
2. CHEMICAL BONDING .............................................................................................. 244
3. ACID AND BASES ....................................................................................................... 256
4. SALTS ............................................................................................................................ 284
5. ELECTRO CHEMISTRY ............................................................................................ 300
6. ANALYTIC CHEMISTRY ........................................................................................... 320
MATHS
1. RELATIONS REVISION AND FUNCTIONS ........................................................... 337
2. REAL NUMBERS AND LIMITS ............................................................................... 348
3. TRIGONOMETRY ...................................................................................................... 362
4. MATRICES ................................................................................................................... 370
5. FUNCTION ................................................................................................................... 380
6. REMAINDER AND FACTOR THEOREM ............................................................... 390
7. QUADRATIC EQUATIONS AND INEQUATIONS .................................................. 404
8. PROGRESSION ........................................................................................................... 426
9. STATEMENTS ............................................................................................................. 438
10. MATHEMATICAL INDUCTION ............................................................................... 456
11. BINOMIAL THEOREM .............................................................................................. 464
PHYSICS
PHYSICS
Contents
1. RIGID BODY DYNAMICS
2. SIMPLE HARMONIC MOTION
3. SOUND
4. MAGNETISM
5. ALTERNATING CURRENTS
6. ELECTRO MAGNETISM
7. THERMO ELECTRICITY
8. SEMICONDUCTOR
9. RAY PPTICS
10. PHYSICAL OPTICS
11. ELECTROLYSIS
1
IIT - FOUNDATION - SET - V
1 RIGID BODY DYNAMICS
2
PHYSICS
Rigid body:
A rigid body is one which does not undergo any change in its shape or size due to the applied force.
Moment of a force:
It is the turning effect of a force about a point. It is equal to product of magnitude of the force and the
perpendicular distance of the point from line of action of force.
If a force ‘F’ acts on a body at point ‘P’ which is at a perpendicular distance ‘d’ from a point ‘O’, then the
moment of the force (or torque) about ‘O’ = F.d.
S.I. unit of torque : newton .metre or N.m
F
O d P
Couple:
A system of two equal, unlike parallel and non-collinear forces constitutes a couple. The effect of a couple
on a body is to produce rotation but not translatory motion.
Moment of couple or torque:
Definition: The turning effect of the couple on the body is called moment of the couple. It is equal to
product of the either force and the perpendicular distance between the lines of action of the two forces.
Moment of the couple = F AB.
Moment of couple has same units and dimensional formula as work or energy.
couple cannot be replaced by a single force.
A B
3
IIT - FOUNDATION - SET - V
A couple can be balanced by another couple in the same plane, having equal and opposite
moment.
Torque is a pseudo vector and it has same units and dimensional formula as work or energy.
The vector relation is r F , hence its magnitude is || r F sin where ‘‘ is
angle between F and r .
Physical significance and definition of moment of inertia:
We know that every body is unable to change its state of rest or uniform velocity, by itself. This
inability is called its inertia in translatory motion. The mass of a body is the measure of its inertia for translatory
motion. In a similar manner, every rotating body is unable to change its uniform angular velocity by itself. This
inability is called its inertia for rotational motion. This is also called rotational inertia or moment of inertia.
Definition: The moment of inertia of a body with respect to any axis is that property by which it is
unable to change its state of rest or of uniform angular motion around that axis.
The ‘moment of inertia’ in rotatory motion plays the same role as ‘mass’ in translatory motion.
Expression for moment of inertia of a body:
r1
m1
r2 m2
Let the body consists of particles of mass m1, m2, .......mn
which are at distances of r1 , r2 ,............rn
from the axis. The moment of inertia I of the body with respect to this axis is given by
I = m1 r12 + m2r22 +...........
n
2
I= m r
i 1
i i
4
PHYSICS
Note: : Thus the moment of inertia of a body depends on (a) the mass of the body and (b) the
axis about which it is rotating or how the matter in the body is distributed around the axis.
: Radius of gyration of the body: Suppose the particles of this body are distributed in a
circle of such radius ‘K’ around the axis that the moment of inertia remains same. Then
I = m1 K2 + m2 K2 + —————
= m i .K 2 = M . K2 ,
I = M.K2
where M represents the mass of the whole body. This radius ‘K’ is called the radius of gyration.
K
Definition : Thus if the moment of inertia of a body of mass M with respect to any axis is I then its radius
of gyration with respect to the same axis is given by
I
K=
M
Rotational K.E of a body:
If a body is rotating with an angular velocity about an axis about which its moment of inertia is ‘I’
1
then its rotational K.E. is given by (K.E)rot = . I 2
2
Let a body be revolving around an axis with an angular velocity ‘’. Each and every particle of their
body has the same angular velocity . If these particles of mass m1, m2 … are at distances r1, r2 … then their
linear velocities are given by
v1 = r1, v2 = r2 ……..
The kinetic energies of these particles are given by
5
IIT - FOUNDATION - SET - V
1
KE1 = ½ m1 v12 = m1 r12 . 2
2
1 1
KE2 = m 2 v 22 m 2 r22 . 2
2 2
……………………………
Rotational KE of the body = Sum of the kinetic energies of the particles.
1 1
KE = m1 r12 . 2 m 2 r22 2 ....
2 2
1
= (m1 r12 m 2 r22 .......) 2
2
1
Rotational KE = .I. 2
2
Where I is the moment of inertia of the body about the given axis.
Theorem of perpendicular axes:
Statement: The sum of moments of inertia of a uniform plane body (lamina) about any two mutually
x p( x,y)
r y
y y
O
X
perpendicular axes in its plane is equal to, moment inertia about another axis perpendicular to its plane and
passing through the point of intersection of the first two axes.
M athematically I x + Iy = Iz
Where
Ix is the moment of inertia of the body about X-axis
Iy is the moment of inertia of the body about Y-axis.
Iz is the moment of inertia of the body about Z-axis.
Where X, Y and Z axes are mutually perpendicular
6
PHYSICS
Proof:
Any particle ‘P’ of mass ‘m’ is at a distance ‘x’ from OY (y-axis) and ‘y’ from OX (x-axis) and
‘r’ from OZ (z-axis).
Then by definition Ix = my2
Iy = mx2
Iz = mr2
= m (x2 + y2) [ from triangle OPN, r2 = x2 +y2 ]
= mx2 + my2
= I y + Ix
Iz = Ix + Iy
Note: This theorem is true only for a lamina (uniform plane body).
Theorem of parallel axes:
Statement: The moment of inertia of a rigid body about any axis is equal to moment inertia about a
parallel axis passing through its centre of mass plus product of mass of the body and square of perpendicular
distance between the two parallel axes.
Mathematically I = Icm + Mr2
y
x
180- 900
O C Q
r
Where I is the moment of inertia about any axis, Icm is the moment of inertia about an axis parallel to first
axis but passing through the centre of mass, ‘M’ is the mass of the body and ‘r’ is the perpendicular distance
between the two parallel axes.
Proof:
Consider a particle of mass ‘m’ at position ‘p’. Let ‘C’ represents centre of mass of rigid body.
Let OP = y, CP = x and OC = r. Let PQ be the perpendicular drawn on to OC extended. Applying cosine rule
in triangle OCP, we have
7
IIT - FOUNDATION - SET - V
y2 = x2 + r2 - 2.x.r. cos(180-)
= x2 + r2 + 2 xr cos.
CQ
= x2 + r2 + 2r. CQ. {from triangle CQP cos = }
x
y2 = x2 + r2 + 2r (CQ).
my2 = mx2 + mr2 + 2r . m. CQ
my2 = mx2 + mr2 + 2r . . m . CQ
But mx2 = Icm = moment of inertia of the body about an axis through its centre of mass.
2
my =I = momnet of inertia of the body about a parallel axis at a distance ‘r’.
mr2 = (m)r2 = Mr2 (since ‘r’ is constant )
m . CQ = 0 ( The algebraic sum of the moments of the masses of the individual
particles about centre of mass is equal to zero).
Now rewriting the equation , we get I = Icm + Mr2 .
Relation between torque ( ) and angular acceleration ():
Consider a particle of mass ‘m’ rotating about an axis at a distance ‘r’ from the axis, under the action of
force ‘F’. Let the linear acceleration of the particle be ‘a’ and its angular acceleration be .
The moment of force or torque = Force perpendicular distance.
F
r
= F. r
= (ma). r
= m. (r) . r [ since a = r ]
= m r2
but mr2 represent the M.I. of the particle about its axis of rotation
=I
i.e. Torque = moment of inertia angular acceleration.
Note:
torque
I=
angular acceleration .
So moment of inertia can also be defined as ratio of torque acting on a body to angular acceleration
produced in it.
8
PHYSICS
Angular momentum (L) :
Definition: Angular momentum of a body is defined as product of its moment of inertia (I) and its
angular velocity () i.e., L = I.
S.I. unit : kg m2s–1 or joule sec.
L2
Relation between (KE) and (L) is KE =
2I
Relation between torque () and angular momentum (L):
Torque () = rate of change of angular momentum.
If the angular momentum changes from L1 to L2 in a time‘t’ seconds, then
L2 L1 d
= (or) = L
t dt
Law of conservation of Angular momentum:
Statement: If no external torque acts on a rotating system, then its total angular momentum is
conserved. i.e., remains constant.
I = constant
i.e In the absence of any external couple, if the moment of Inertia of the rotating body change from I1 to
I2 due to redistribution of matter then there will be a corresponding change in angular velocity 1 to 2
such that I11 = I2 2.
Illustrations:
A person diving into a swimming pool draws his hands and legs nearer to the trunk. By doing so,
he decreases moment of inertia of the body, thereby increasing his angular velocity. This enables him
to make somersaults.
In the similar way ballet dancers, ice - skaters can increase or decrease the angular speed of spin
about a vertical axis.
To find expressions for the M.I. for different regular bodies about certain
standard axes:
1) To find the moment of inertia of a thin circular ring about its natural axis.
Let M be the mass of the ring, R be the radius of the ring and m be the mass per unit length of the ring.
M = m. 2R
Consider an elementary portion of length d . Its mass is m.d. The M.I. of this elementary portion
about the natural axis of the ring is given by
dI=(m.d)R2 .
The M.I. of the whole ring is given by I = dI
9
IIT - FOUNDATION - SET - V
m.d.R mR 2 d
2
=
= mR2 () = mR2 . 2R
I = MR2
O R
d1
2) To find the M.I. of a circular disc or a solid cylinder about its natural axis.
Let M = Mass of the disc
R = Radius of the disc
m = mass per unit area of the disc
M = m. R2
A thin ring of radius ‘r’ and of thickness ‘dr’ of the
disc is considered as elements portion.
area of the ring = 2r dr
mass of the ring = m.2r.dr
The M.I. of this elementary portion about the natural axis of the disc is given by
dI = (m. 2r.dr). r2 = m. 2 r3 dr
The M.I. of the whole disc is given by
I = dI = m.2. r 3 . dr
R
= m. 2 r 3 . dr
0
10
PHYSICS
R 4 m. R 2 . R 2 MR 2
= m. 2
4 2 2
MR 2
I=
2
Extending this theory to a solid cylinder, the M.I. of
a cylinder about its natural axis is also given by the same
MR 2
formula, I=
2
The moment of inertia of the cylinder about an axis through its centre and perpendicular to its
2 R2
natural axis is given by I = 12 4 , where ‘ ’ is length of the cylinder and ‘R’ is its radius.
M
For a hollow cylinder the M.I. about its natural axis is given by I = MR2.
3) To find the M.I. of a thin uniform bar about an axis through its mid point and
perpendicular to its length.
Let M = mass of the bar and L = length of the bar
Let m = mass per unit length of the bar
M = m.L
Consider an elementary portion of length ‘d’ at a distance ‘ ’ from the axis.
Mass of elementary portion = m. d
The M.I. of the elementary portion about the given axis is
dI = m.d . 2
The M.I. of the whole bar is given by
L
I = dI = m.d. 2 = m 2L 2 .d
2
11
IIT - FOUNDATION - SET - V
A B
L L
2 2
m. L3 mL. L2 ML2
=
12 12 12
ML2
I=
12
From the theorem of parallel axes the moment of inertia of the bar about an axis through its one
end and perpendicular to its length is given by
I1 I c Md 2
M 2
M
=
12 2
1M 2
I
3
4) Moment of inertia of a rectangular lamina
Mass of the lamina = M
Length of the lamina = 1
Breadth of the lamina = b
Let the lamina be divided into large number of thin strips Parallel to its length. Each strip of mass ‘dM’
acts as a thin Bar of length ‘’, moment of inertia of the strip about an axis AB along the edge of the rectangular
lamina as shown in the
Figure is given by
dM
12
PHYSICS
2
dI = dM
3
The M.I of the whole rectangular lamina about the axis AB is given by
2 M. 2
I = dI dM.
3 3
From the theorem of parallel axis, the moment of inertia of the Lamina (Iy) about an axis YY’ through
the centre of mass and parallel to AB is given by
A
Y
d
2
CM
B Y1
I = Iy + Md2
2
M 2
I y M
3 2
M 2 M 2
Iy
3 4
M 2
Iy ……….. (1)
12
Similarly the moment of inertia of the lamina about an axis XX’ through the centre of mass and parallel to
length is given by
X1 X
CM
Mb 2
Ix
12
From the theorem of perpendicular axis, the moment of inertia of the lamina about Z-axis through its
centre of mass and perpendicular to its plane is given by
Iz Ix I y
Mb 2 M 2
Iz
12 12
13
IIT - FOUNDATION - SET - V
2 b2
I
2 M
12 12
Z Y
X1 X
Y1
14
PHYSICS
ASSIGNMENT
15
IIT - FOUNDATION - SET - V
1. Among wooden sphere, copper sphere of equal mass which one has greater moment of inertia about an
axis passing through their diameter.
1) copper sphere 2) wooden sphere
3) both have same moment of inertia 4) None
2. AB is a stick half of which is wood and the other half is steel. IA is moment of inertia about an axis
passing through A and perpendicular to its length, IB is the moment of inertia about an axis passing
through B and perpendicular to its length. Then
1) IA = IB 2) IA > IB 3) IB > IA 4) none
Wood Steel
A B
3. Moment of inertia of a solid sphere about its diameter is I. If that sphere is recast into 8 identical small
spheres, then the moment of inertia of such small sphere about its diameter is
1) I/8 2) I/16 3) I/24 4) I/32
4. Three solid spheres each of mass 10 kg and diameter 0.2m are placed at the three corners A, B, C of
an equilateral triangle of side 1m. The M.I. of the system about AB axis is
1) 0.12 kg-m2 2) 7.5 kg-m2 3) 7.62 kg-m2 4) None
5. A thin wire of mass m and length L is bent in the form of a ring. Moment of inertia of the ring about an axis
passing through its centre and perpendicular to its plane is
mL2 mL2 mL2 mL2
1) 2) 3) 4)
4 2 4 2 2 2
6. Ratio of densities of materials of two circular discs of same mass and thickness is 2 : 3. Ratio of their
moment of inertia about their natural axes is
1) 2 : 3 2) 3 : 2 3) 4 : 9 4) 1 : 1
7. The masses 1kg, 2kg and 3kg are placed at A, B, C respectively, and are joined with thin rods of
negligible mass. If AB=3m, BC=4m and CA=5m then their moment of inertia about an axis passing
through A and perpendicular to the plane ABC is
1) 93 kg m2 2) 125 kg m2 3) 107 kg m2 4) None
8. A straight thin uniform rod of length ‘4L’ and mass ‘4M’ is bent into a square. Its M.I. about one side is
5ML2 7 ML2 ML2 ML2
1) 2) 3) 4)
3 6 48 3
9. Moment of inertia of a thin square plate about an axis passing through its centre and perpendicular to its
plane is 100 kg m2. Its moment of inertia about an axis passing through its diagonal is
1) 400 kg m2 2) 200 kg m2 3) 100 kg m2 4) 50 kg m2
16
PHYSICS
10. Moment of inertia of a thin square plate about an axis passing through its diagonal is I. Its moment of
inertia about an axis passing through its centre in the plane of the plate and making an angle with the
diagonal is
1) I 2) I cosq 3) I sinq 4) I tanq
11. M.I of a thin circular ring of radius R and mass M about a chord at a distance equal to half of the radius
of the ring is
5 2 MR 2 3 3
1) MR 2) 3) MR 2 4) MR 2
4 4 2 4
12. Four spheres each of mass M and diameter 2r, are placed with their centres on the four corners of a
square of side a (>2r). The moment of inertia of the system about one side of square is
2 2
1)
5
M ( 5r 2 4 a 2 ) 2)
5
M 5r 2 2 a 2
2
3) 2M (2r2+5a2)/5 4)
5
M 4 r 2 5a 2
13. Two loops P and Q are made from a uniform wire. The radii of P and Q are r1 and r2 respectively and
r2
their moments of inertia are I1 and I2 respectively. If I2/I1=4 then equals
r1
1) 42/3 2) 41/3 3) 4-2/3 4) 4-1/3
14. The radius of a circular disc is 7cm. The radius of gyration of the disc for rotation about its natural axis
is nearly equal to
1) 3.5 cm 2) 14cm 3) 1.414 cm 4) 5 cm
15. A mass M is moving with a constant velocity parallel to the x-axis. Its angular momentum with respect
to the origin
1) zero 2) remains constant
3) goes on increasing 4) goes on decreasing
16. If all the persons on the surface of the earth goes to Antarctica then the duration of the day
1) increases 2) decreases 3) remains unchanged 4) None
17. Mass remaining same, if the earth shrinks until its radius becomes half of its value then the length of the
day will be
1) 8 hrs 2) 6 hrs 3) 12 hrs 4) 16 hrs
18. A thin circular ring of mass M and radius R is rotating about its axis with a constant angular velocity w.
Two objects each of mass m are attached gently to the ring. The ring now rotates with an angular
velocity
M 2 m M M 2 m
1) w M(M+m) 2) 3) 4)
M 2 m M 2 m M
17
IIT - FOUNDATION - SET - V
19. A circular disc is rotating about its natural axis with an angular velocity of 10 rad/sec. A second disc of
same mass is joined to it coaxially. If the radius of this disc is half of the radius of the first then they
together rotate with an angular velocity of
1) 2.5 rad/s 2) 5 rad/sec. 3) 8 rad/s 4) 6.67 rad/s
20. A ballet dancer spins about a vertical axis at 60 r.p.m with her arms closed. If she now stretches her
hands M.I increases by 50%. Her new speed of revolution is
1) 80 r.p.m 2) 40 r.p.m 3) 90 r.p.m 4) 30 r.p.m
21. A solid sphere and a solid cylinder are rolling down without slipping on an incline plane. The ratio of
their accelerations is
1) 1 : 1 2) 9 : 10 3) 15 : 14 4) 10 : 9
22. What fraction of the total energy of a rolling circular disc is its translatory K.E.
1) 2/3 2) 1/2 3) 3/4 4) 1/4
23. A circular disc and a ring of same mass and radius are rolling with same linear velocity. Their K.E. are
in the ratio of
1) 1 : 1 2) 4 : 3 3) 3 : 4 4) 3 : 2
24. A cylinder of mass 10kg and of radius 0.1 m is rolling with a linear velocity of 2 m/s. The Kinetic
Energy of the rolling cylinder is
1) 28 joule 2) 40 joule 3) 30 joule 4) 80 joule
25. A metre scale is held vertically on the table. It is allowed to fall down without shipping at its bottom. The
linear velocity with which the top end hits the table is
3g g
1) 3g 2) 3) g 4)
2 2
26. A uniform rod PQ of length L is hinged at one end P. The rod is kept in the horizontal position by a
mass-less string tied to point Q. If the string is cut, the angular velocity of the rod after an angular
displacement of 300 is
g 2g 3g 3g
1) 2) 3) 4)
L L 2L L
String
P Q
<-----------L----------------->
18
PHYSICS
27. A uniform rod PQ of length L is hinged at one end P. The rod is kept in the horizontal position by a
mass-less string tied to point Q. If the string is cut, the initial angular acceleration of the rod will be
1) g/L 2) 2g/L 3) 6g/L 4) 3g/2L
String
P Q
<-----------L----------------->
28. Two bodies A, B have M.Is Ia, Ib (Ia > Ib) have the same angular momentum. Then kinetic energy is
more for
1) body A 2) body B
3) both have same kinetic energy 4) None
29. A cylinder has the same M.I about an equatorial axis and also about its own axis. Then the ratio
between its length and radius is
1) 3 : 1 2) 1 : 3 3) 2 : 1 4) 1 : 2
30. The radius of gyration of a body about an axis at a distance of 4cm from the centre of gravity is 5cm. Its
radius of gyration about a parallel axis through the centre of gravity is
1) 31 cm 2) 3cm 3) 1cm 4) 2cm
31. Two point masses 3kg and 9kg are joined by a rod of length 1m. and of negligible mass. Moment of
Inertia of the system about an axis passing through its centre of mass and normal to the rod is
1) 6 kg m2 2) 2.25 kg m2 3) 3 kg m2 4) None
32. A particle of mass 4kg is projected with a velocity 2 m/s making an angle of 450 with the horizontal.
The magnitude of the angular momentum of the particle about the point of projection when the particle
is at the maximum height of its path is
( g = 10m/s2)
1) 0.2 kgm2/s 2) zero 3) 5kg m2/s 4) 2.5 kgm2/s
33. Moment of inertia of a uniform circular disc about a diameter is I. Its moment of inertia about an axis
perpendicular to its plane and passing through a point on its rim will be
1) 2I 2) 3I 3) 5I 4) 6I
34. A solid cylinder of mass ‘m’ rolls without slipping down an inclined plane making an angle q with the
horizontal. The frictional force between the cylinder and the inclined plane is
1) mg sinq 2) (mg sinq ) / 3 3) mg cosq 4) (2mg sinq ) / 3
35. Reel of thread in the form of solid cylinder is allowed to unroll by holding the loose end of thread in hand.
The acceleration with which the reel falls down is
(g = 9.9 m/s2)
19
IIT - FOUNDATION - SET - V
1) 9.9 m/s2 2) 3.3 m/s2 3) 6.6 m/s2 4) none of the above
36. A tangential force of 10N is acting on a circular plate of radius 50cm such that it can rotate about an axis
perpendicular to its plane passing through its center. If its moment of inertia about the given axis is
0.5kgm2, number of revolutions it makes in the first 6 seconds is
90 180
1) 2) 3)90 4)180
37. A solid sphere and a hollow sphere of same mass and radii are rolling down an inclined plane from the
same point. The ratio of their translational kinetic energy at the bottom of the plane is
1) 1 : 1 2) 25 : 21 3) 7 : 5 4) 5 : 3
38. A body is rotating with an angular velocity 4rad/s and its radius of gyration is 2m. If the angular velocity
becomes four times without any external torque then the radius of gyration will become
1) 4m 2) 0.5m 3) 1m 4) 8m
39. A metal rod has moment of inertia I about an axis passing through its center and perpendicular to its
length. It is bent at the middle such that two parts are perpendicular to each other, and perpendicular to
the axis. The moment of inertia of the system about the same axis will be
1) 2I 2) I 3) I/2 4) 4I
40. The kinetic energy of body rotating at a constant rate of 5rev/s is 62.8J. Its angular momentum is
1) 12kgm2/s 2) 2.4kgm2/s 3) 3.8kgm2/s 4) 4kgm2/s
41. A simple pendulum of mass m and length L is vibrating with angular amplitude 600 The torque acting on
the bob at the extreme position with respect to point of suspension is
mgL 3mgL
1) 2) mgL 3) 4) zero
2 2
42. Three particles of masses 3kg, 2kg and 5kg are situated at (3m, 0,0)(0, 2m, 0), (0,0,1m) respectively.
Radius of gyration of the system of particles about the Z-axis is
13
1) 2m 2) m 3) 2.82m 4) 3.5 m
6
43. Four identical thin rods, each of mass ‘M’ and length ‘L’ are joined to form a rigid square frame. The
frame lies in the xy plane with its centre at the origin and the sides parallel to the x and y axes. Its moment
of inertia about an axis parallel to the Z-axis and passing through one corner is
4ML2 10ML2
1) 2) 3) 4ML2 4) 3ML2
3 3
44. An inclined plane of length 3m is held at an angle 300 with the horizontal. A hollow sphere rolls from the
top without slipping. The linear velocity of the center of mass with which it arrives the ground is
1) 1m/s 2) 4.2m/s 3) 2.4m/s 4) 2.1m/s
20
PHYSICS
45. F = ai + 3j + 6k ; r = 2i - 6j -12k. The value of ‘a’ for which the angular momentum is conserved, is
1) -1 2) 0 3) 1 4) 2
46. A man is spinning in the gravity free-space changes the shape of the body by spreading his arms. By
doing this he can change his (a) moment of inertia
(b) angular momentum (c) angular velocity (d) rotational kinetic energy
Which one of the following are correct
1) a, b and c 2)d, a and b 3) c, d and a 4) b, c and d
47. Assertion(A): Two solid cylinders of different masses and radii take same time to roll down an inclined
plane without slipping
Reason(R): Moment of inertia of two solid cylinders of different masses and radii about their natural
axes are different.
1) Both A and R are true and R is the correct explanation of A.
2) Both A and R are true but R is not the correct explanation of A.
3) A is true but R is false.
4) A is false but R is true.
48. Statement - A: If a uniform metal disc is remoulded into a solid sphere, then the moment of inertia about
the axis of symmetry increases than that before.
Statement - B: For a given body and for a given plane, the moment of inertia is minimum about an axis
passing through the centre of mass.
1) Both A and B are wrong. 2) Both A and B are correct.
3) A is correct and B is wrong. 4) A is wrong but B is correct.
49. Consider the following two statements.
A: Linear momentum of the system remains constant.
B: Centre of mass of the system remains at rest.
1) A implies B and B implies A 2) A does not imply B and B does not imply A.
3) A implies B but B does not imply A 4) B implies A but A does not imply B.
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IIT - FOUNDATION - SET - V
KEY
1. 2 16. 2 31. 2 46. 3
2. 2 17. 2 32. 1 47. 2
3. 4 18. 3 33. 4 48. 2
4. 3 19. 3 34. 2 49. 4
5. 3 20. 2 35. 3
6. 2 21. 3 36. 1
7. 1 22. 1 37. 2
8. 1 23. 3 38. 3
9. 4 24. 3 39. 2
10. 1 25. 1 40. 4
11. 4 26. 3 41. 3
12. 4 27. 4 42. 1
13. 2 28. 2 43. 2
14. 4 29. 1 44. 2
15. 2 30. 2 45. 1
22
PHYSICS
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IIT - FOUNDATION - SET - V
SIMPLE HARMONIC
2 MOTION
24
PHYSICS
Definition of SHM:
The motion of a particle is said to be simple harmonic, if its acceleration is always directed
towards a fixed point in its path and is directly proportional to its displacement from that point.
Explanation:
O a
x
Let a particle be executing S.H.M. with ‘O’ as its mean position and let ‘x’ be the displacement at an
instant. If the acceleration at this instant is ‘a’ then,
by definition, a - x
or the force F -x
F = - k x, where k is called as the force constant.
S.I. unit of Force constant : Nm-1
To show that the projection of uniform circular motion on any diameter is simple harmonic:
Consider a particle ‘P’ in uniform circular motion. Let ‘r’ be the radius of circle and ‘‘ be its
constant angular velocity.
Let the angular displacement of the particle be ‘‘ in ‘t’ sec.
= t
N r 2 cos P
r2sin
y r
X1 O M X
Y1
As the particle ‘P’ moves along the circumference of the circle its projection ‘N’ that is the foot of the
perpendicular ‘N’ executes SHM along YY1. As the particle moves from X to P, its projection on YY1 moves
from O to N.
Displacement of the projection y = ON
= OP sin
y = r sin t ( = t)
2
Then centripetal acceleration of the particle ‘P’ is ‘r ’ and it acts in the direction P to O.
This acceleration can be resolved into two rectangular components:
r2 cos along PN and r2 sin along PM.
Acceleration of the projection ‘N’ is given by the component r2sin which is parallel to YY1.
a = r 2 sin = 2 r sin
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IIT - FOUNDATION - SET - V
a = 2 (-y )
Negative sign indicates ‘y’ and ‘a’ are opposite in direction.
a - y ( is constant)
The motion of the projection ‘N’ along YY1 is SHM.
Time period of a body in SHM : Time period is the time taken by a particle to complete one oscillation.
As the particle ‘P’ completes one rotation on the circumference of the reference circle, the foot
of the perpendicular ‘N’ completes one oscillation. Therefore, the time periods of both are equal.
2
So, T = but a = 2y
2 y displaceme nt
T= 2 T 2
a/y a accelerati on
The above expression is the time period of a particle in SHM in terms of acceleration and
displacement of a particle.
Phase:
The phase of a particle in S.H.M. at any instant represents the state of its vibration at that instant. In other
words it gives the details regarding its position and direction of motion.
N P
+
S
y = t
X1 O X
Y1
If the particle passes from the initial position S to P in ‘t’ sec, then
y = r sin ( + )
= r sin (t + )
The angle (t + ) gives the phase at the instant ‘t’.
At the instant t=0
y = r sin [ = t = 0 ]
So ‘‘ is called initial phase or phase constant or epoch.
26
PHYSICS
The displacement y in general is given by
y = r sin(t )
Note: Angle ‘ ‘ is also called the phase difference between the positions P and S.
Properties of a particle in SHM:
While the body is in uniform circular motion with angular velocity along a circle of radius r, its
projection N is in SHM along YY1.
Linear velocity of the body at any instant is V = r and acts along the tangent.
Acceleration of a body is a = r 2 and acts along the radius.
a) Displacement of N, y = rsint
where r is the amplitude of vibration
b) Velocity of N, V = component of the
velocity of the body parallel to YY1.
= Vcos = rcost
= r 1 sin 2 t
y2
= r = 1 = r 2 y2
r2
Velocity of N, V = r cost
or V=
V=r Vcos
Y
Vsin
N P
r2 r2sin
X1 O M X
Y1
c) Acceleration of N = a = component of the acceleration of the body parallel to YY1
= r 2sin
= 2y
Acceleration of N is given by a = 2y
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IIT - FOUNDATION - SET - V
a = 2r sint
d) KE of the particle in SHM :
If the projection ‘N’ is treated as a particle of mass ‘m’ then,
1
its K.E. = mv 2
2
but its velocity v = r . cost = r 2 y 2
1 22 2
K.E = m r cos t
2
And the KE at the instant when the displacement is ‘y’ is given by
1
or K.E = m 2 (r2 - y2)
2
Y
K.E.
X X
A O B
displacement
The above graph shows the variation of K.E. of the particle with its displacement from
the mean position A and at the extreme positions.
OA = OB = amplitude.
e) The total energy of the particle in SHM :
Velocity of the particle in SHM, v = r 2 y2
1
KE of the particle in SHM, KE = m 2 ( r 2 y 2 )
2
1 1
KE of the particle in Mean position = m 2 (r 2 0) m 2 .r 2
2 2
1
Total energy = KE in the Mean position = m2 r 2
2
1
T.E. = m2r2
2
f) P.E. of the particle in SHM :
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PHYSICS
Applying the law of conservation of energy
P.E. = T.E. - Kinetic Energy
1 1
= m2r2 - m2 (r2 - y2)
2 2
1
= m2y2
2
P.E. = m2 y2
1
or P.E. = m2 r2 sin2 t
2
Y
P.E.
X X
A O B
displacement
The variation of P.E. of the particle with its displacement from the mean position is
shown in the graph. P.E. is maximum at the extreme position A and B P.E. is minimum at the mean position.
2
2 acceleration
g) Period of oscillation of N = T = = ²
displacement
1
h) Frequency of N = n = ³
T 2
Simple pendulum:
A simple pendulum consists of a small bob, suspended by a light inextensible thread.
Expression for time period:
Let ‘S’ be the point of suspension of simple pendulum. Let ‘’ be its length, which is equal to
distance between point of suspension and the centre of mass of the bob.
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IIT - FOUNDATION - SET - V
T
A B
mgsin
mgcos
mg
Let ‘A’ be its rest or mean position. Let the bob of simple pendulum be displaced through a
small angle (less than 50) and released. Let the bob be in the position B at some instant, when the
angular displacement is ‘’.
Then AB = x = displacement from mean position.
At position ‘B’ the forces acting on the bob are the weight mg and the tension ‘T’ in the string.
Weight ‘mg’ of the bob is acting vertically downwards. This weight ‘mg’ can be resolved into two
rectangular components ‘mgsin‘ and ‘mgcos‘ as shown in the figure.
Tension ‘T’ in the string balances ‘mgcos’. The component ‘mgsin’ pulls the bob back towards
the mean position.
Restoring force on the bob is F = - mg sin
restoring force mg sin
acceleration =- = - g sin
mass m
= - g . () { if ‘‘ is small, then sin }
x arc AB x
= - g. since
radius r
g
a = . x
a - x { since ‘g’ is constant at a given place and ‘’ is constant for a given pendulum }.
So the motion of bob of the simple pendulum is SHM.
But for a body in SHM, a = - 2 x
g g
From equations & we get, 2 = =
2 2
But we know T = g or T = 2
g
Where ‘T’ is time period of simple pendulum and ‘g’ is acceleration due to gravity at a given place.
30
PHYSICS
A simple pendulum for which time period is 2 seconds is called seconds pendulum. Its length is
approximately 1 metre on the surface of the earth.
Laws of Simple Pendulum:
Time period of simple pendulum is independent of mass, size and material of the bob.
Time period is independent of amplitude provided it is less than 40.
At a given place time period of simple pendulum is directly proportional to square root of its
length, i.e. T .
The time period of a given simple pendulum is inversely proportional to square root of the accel-
eration due to gravity at the place of the experiment.
1
T g
Oscillations of a loaded spring:
When a load of mass M is suspended from a spring then it is stretched by an amount ‘x’ and
comes to rest. The tension ‘T’ in the stretched spring balances the weight Mg.
Stretching force = constant elongation
Mg = k.x
The constant ‘k’ is called the force constant of the spring or simply the spring constant.
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IIT - FOUNDATION - SET - V
If the load is pulled down from its mean position and released then it begins to oscillate about its
mean position.
At the instant when the displacement of the load from its mean position is ‘y’ the additional
tension developed in the spring= ky.
This tension accelerates the body towards its mean position.
k. y k
Acceleration = = . y
M M
Acceleration = constant displacement but the
constant is 2
k k
2 = or =
M M
This load executes S.H.M. with a period given by
2
T=
2
T=
M
T = 2
k
x
From equation it also follow that T = 2
g
To show that the law of conservation of energy holds good for Simple Pendulum:
Let a bob of mass ‘m’ be executing S.H.M with an amplitude ‘a’. The displacement ‘y’ of this bob at any
instant ‘t’ is given by
y = a.sint.
Where is a constant
The velocity at this instant is given by
v = a cost
The acceleration at the same instant is given by
A = -2.asint = -2y
So retardation =2y and retarding force = m2y
Work done by the retarding force for any small displacement ‘dy’ is given by
dw = m2y.dy
32
PHYSICS
t
P
O
Total work done during its displacement from O to P is given by
y
2 1 2 2
W= m y.dy 2 m y
0
1
W= m2a2sin2t
2
This work is stored in the bob as P.E.
1 1
P.E. in the position P = m2a2sin2t = m2y2
(1)
2 2
1 1
Also, K.E. in the position P = mv2 = m2a2cos2t
(2)
2 2
Total energy at P = P.E. + K.E.
1 1
T.E. = m2a2sin2t + m2a2cos2t
2 2
1
T.E. = m2a2
(3)
2
Thus the total energy is independent of time ‘t’ and displacement ‘y’. So the T.E remains constant
throughout. Hence the law of conservation of energy holds good.
33
IIT - FOUNDATION - SET - V
ASSIGNMENT
34
PHYSICS
1. The displacement of a particle in S.H.M. is x = 3 sin (20pt) + 4 cos(20pt) cm . Its amplitude of
oscillation is
1) 3 cm 2) 4 cm 3) 5 cm 4) 25 cm
2. A particle oscillates as per the equation x = (7cos0.5pt) m , the time taken by the particle to move from
the mean position to a point 3.5m away is
1) 1/3 s 2) 1/2 s 3) 1 s 4) 2/3 s
3. Two particles are in S.H.M along parallel straight lines with same amplitude and time period. If they
cross each other in opposite directions at the mid point of mean and extreme position, phase difference
between them is
1) 300 2) 1200 3) 1500 4) 1800
4. The period of a particle in SHM is 8 seconds. At t=0, it is at the mean position. The ratio of the distances
travelled by it in the first and the 2nd second is
1) 1/2 2) 1/ 2 3) 2 4) +1
5. Two particles are in SHM with same amplitude and frequency along the same line and about the same
point. If the maximum separation between them is equal to their amplitude, the phase difference be-
tween them is
1) p/2 rad 2) /3 rad 3) p/6 rad 4) 2p/3 rad
6. The period of oscillation of a particle in SHM is 4s and the amplitude of vibration is 7 cm. The velocity
of the particle 2/3s after passing the mean position is ........
1) 5 cm/s 2) 5.5 cm/s 3) 6 cm/s 4) 7 cm/s
7. The velocity of a particle in SHM at the instant when it is 0.6cm away from the mean position is 4 cm/
s. If the amplitude of vibration is 1 cm then its velocity at the instant when it is 0.8 cm away from the
mean position is
1) 2.25 cm/s 2) 2.5 cm/s 3) 3.0 cm/s 4) 3.5 cm/s
8. A simple pendulum performs simple harmonic motion about x=0 with an amplitude ‘A’ and time period
‘T’. Speed of the pendulum at x=A/2 will be
3
1) A 3 / T 2) A / T 3) A 4) 3 2 A / T
2T
9. A hole is bored along the diameter of the earth and a stone is dropped into it. If the radius of the earth is
R, then the speed of the stone, when it is at the centre of the earth, is
gR
1) 2) gR 3) 2gR 4) Zero
2
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IIT - FOUNDATION - SET - V
10. The period of oscillation of a particle in SHM is p second and its amplitude of vibration is 10 cm. The
acceleration of the particle p/12 second after passing the mean position is
1) 20 cm/s2 2) 20 cm/s2 3) 20/ cm/s2 4) 20 cm/s2
11. The velocity of a particle in SHM at the mean position is 1ms-1 and acceleration at the extremity is 2ms-
2
. The angular frequency of the particle is
1) 2 rad s-1 2) 1 rad s-1 3) 0.5 rad s-1 4) 3 rad s-1
12. If the displacement ‘x’ and velocity ‘v’ of a particle in SHM along a straight line are connected by the
relation 4v2 = 25 - x2, then its time period is
1) 2p second 2) 5p second 3) 4p second 4) 2.5p second
13. A horizontal board is made to perform SHM along a horizontal straight line of length 32 cm.A body is
placed on the table (m=0.2). The maximum number of oscillations per second that can be made by the
board so that the body will not slip on the board is
1) 7/2 p 2) 7 / 4 p 3) 7 / 8 p 4) None
14. A horizontal platform is executing SHM up and down with a period of 1s. The maximum amplitude with
which it can vibrate so that an object placed on it does not leave it is (p2 = g)
1) 0.25m 2) 0.5m 3) 1m 4) 1.25m
15. The displacement of a particle in SHM whose amplitude is A at which PE is ¼ th of total energy is
A A A A
1) 2) 3) 4)
2 2 4 2 2
16. The amplitude of a particle in SHM is 0.2m. The P.E and K.E of this particle 0.1sec after passing the
mean position are equal. The period of oscillation is
1) 1 sec 2) 0.8 sec 3) 0.6 sec 4) 0.4 sec
17. An object of mass 0.2 kg is in SHM along X- axis with a frequency of hertz. At the position x=0.04 m,
it has KE of 0.5J and PE of 0.4 J. The amplitude of vibration is
1) 0.06 m 2) 0.05m 3) 0.08m 4) 0.09m
18. A linear harmonic oscillator of force constant 2 ´ 106 N/m and amplitude 0.01m has a total mechanical
energy of 160J . Its
1) maximum potential energy is 100 J 2) maximum kinetic energy is 100 J
3) minimum potential energy is zero 4) maximum kinetic energy is 160 J
19. If the length of the simple pendulum increases by 44% the time period increases by
1) 12% 2) 20% 3) 44% 4) 40%
20. The length of a seconds pendulum on earth is 1m. What would be the length of such a seconds pendu-
lum on a planet whose mass is the same as that of earth and whose radius is twice the radius of the earth?
1) 0.25m 2) 4m 3) 0.5m 4) 2m
36
PHYSICS
21. A simple pendulum with a brass bob has a period T. The bob is now immersed in a non-viscous liquid
and oscillated. If the density of the liquid is th of brass, the time period of the same pendulum will be
8 8 64
1) T 2) T 3) T 4) T
7 7 49
22. Two pendulums of lengths 1.21 m and 1m are in phase at a given instant of time. The minimum number
of oscillation completed by the shorter pendulum after which again they will be in phase is
1) 5 2) 6 3) 10 4) 11
23. Two SHMs have time periods T and . They start SHM in the same phase. After the larger pendulum
performs one oscillation the phase difference between them is
1) 450 2) 900 3) 600 4) 300
24. A clock S is based on oscillations of a spring and a clock P is based on pendulum motion. Both clocks
run at the same rate on Earth. On a planet having the same density as Earth, but twice the radius
1) S will run faster than P 2) P will fun faster than S
3) They will both run at the same rate as on Earth 4) None of the above
25. A spring of force constant K, is cut into two parts, such that their lengths are in the ratio of 2 : 3. The
force constant of the shorter part is
1) 5 K 2) 2.5K 3) 0.4 K 4) 1.5K
26. When a mass is freely hung from two springs separately the period of vertical oscillations are T1 and T2.
When the mass is hung from the same two springs connected in series, the period will be
T1 T2 T12 T22
1) 2) 3) T12 T22 4) (T12 T22 )
2 2
27. When a body is suspended from two light springs separately, the periods of vertical oscillations are T1
and T2 . When the same body is suspended from the two springs connected in parallel, the time period
will be
2T1T2 T1T2
1) T1 2 T2 2 2) T1T2 3) 4)
2 2
T1 T2 T1 2 T2 2
28. A mass M is suspended from a spring of negligible mass. The spring is pulled a little and then released so
that the mass executes simple harmonic motion with time period T. If the mass is increased by ‘m’ then
the time period becomes . The ratio is
1) 9/16 2) 25/16 3) 4/5 4) 5/4
29. Two bodies M and N of equal masses are suspended from two separate massless springs of spring
constants K1 and K2 respectively. If the two bodies oscillate vertically such that their maximum veloci-
ties are equal, the ratio of the amplitude of M to that of N is
1) K1 : K2 2) K 1 : K2 3) K2 : K1 4) K2 : K1
37
IIT - FOUNDATION - SET - V
30. If there is a planet on which the value of the acceleration due to gravity is one-nineth of that on the earth’s
surface the frequency of oscillation of a simple pendulum on the planet will be
1) one-third 2) one-ninth 3) nine times 4) none
31. The time-period of vertical oscillations of a load attached to the parallel combination of two identical
springs is 2s. If the same load is attached to the series combination of the same identical springs, then
the time period of the oscillating load is
1) 1s 2) 2s 3) 3s 4) 4s
32. A mass M attached to a spring oscillates with a period of 2 seconds . If the mass is increased by 2kg
, then the period increases by 1second . The initial mass is
1) 1.6 kg 2) 4kg 3) 2kg 4) 2pkg
33. Pendulums of lengths and have the same time-period T at two different places. If = x, the difference is
the value of ‘g’ at those two places is
4x 4x 4 2 x
1) 2) 3) 4) None
T T2 T2
34. A pendulum clock is beating seconds at a place where acceleration due to gravity is 980 cms-2. If it is to
show correct time at a place where acceleration due to gravity is 982 cms-2, its length should be
1) increased by cm 2) increased by cm
3) decreased by cm 4) decreased by cm
35. A simple pendulum is beating seconds in a stationary lift. If the lift is moving up with an acceleration g/4
then its period will be
1) 2seconds 2) seconds 3) seconds 4) seconds
36. A mass of 1kg stretches a spring by 0.8m. If the mass is further pulled down by 0.05m and released,
then the mass oscillates with a period (g = 9.8ms-2)
1) p seconds 2) 2pseconds 3) seconds 4) seconds
37. A particle executing SHM had a maximum velocity of 40cms-1 and maximum acceleration of 80cms-2.
Then its amplitude and the period of oscillation are
1) 20cm, p/2sec 2) 10cm, p/2sec 3) 20cm, psec 4) 10cm, p sec
38. A particle performing SHM has velocities of 3cms-1 and 4cms-1 respectively when the displacements are
4cm and 3cm. The maximum velocity of the particle is
1) 5cms-1‘ 2) 6cms-1 3) 7cms-1 4) 8cms-1
39. A particle is oscillating with an amplitude of 12cm and time period 12second. The shortest time in
which the particle moves from the position +6.0cm to –6.0cm is
1) 2 second 2) 4 second 3) 8 second 4) 12 second
38
PHYSICS
40. A particle executes simple harmonic motion with frequency f. The frequency with which its kinetic
energy varies is
1) f/2 2) f 3) 2f 4) 4f
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IIT - FOUNDATION - SET - V
KEY
1. 3 16. 2 31. 4
2. 1 17. 1 32. 1
3. 2 18. 2 33. 3
4. 4 19. 2 34. 2
5. 2 20. 1 35. 3
6. 2 21. 1 36. 3
7. 3 22. 4 37. 3
8. 1 23. 2 38. 1
9. 2 24. 2 39. 1
10. 2 25. 2 40. 3
11. 1 26. 3
12. 3 27. 3
13. 2 28. 1
14. 2 29. 4
15. 2 30. 1
40
PHYSICS
41
IIT - FOUNDATION - SET - V
3 SOUND
42
PHYSICS
Audible Sounds: Sounds produced by a vibrating body whose frequency lies between 20Hz and 20,000Hz
are called audible sounds.
Inaudible Sounds:
a) Infrasonics:Sounds of frequency less than 20Hz are called infrasonics”.
Ex: Sound from oscillating simple pendulum.
b) Ultrasonics : Sounds of frequency greater than 20,000Hz are called
“ultrasonics”.
Ex: Sound produced by bats.
Both infrasonics and ultrasonics are inaudible to human ear.
Propagation of Sound:
A source of sound disturbs the surrounding medium. These disturbances namely the compressions
and rarefactions travel in the surrounding medium in the form of a wave. Thus, sound wave is basically a
mechanical wave and hence a material medium is necessary for the propagation of sound wave.
There are two types of waves namely, Transverse wave ‚ Longitudinal wave.
Differences between Transverse and Longitudinal waves:
Transverse Wave Longitudinal Wave
1 The particles of the medium vibrate at 1 The particles of the medium vibrate
right angles to the direction of along the direction of propagation of
propagation of the wave. the wave.
2 A crest and a trough make up a 2 A compression and a rarefaction make
complete wave. up a complete wave.
3 It can travel in solids only. 3 It can travel in solids, liquids and
gases.
4 It exhibits polarisation. 4 It does not exhibit polarisation.
Ex: Vibrations of stretched strings. Ex: Sound waves.
Representation of wave:
A transverse or a longitudinal wave can be represented by a sine curve. Such a sine curve is obtained
when a graph is drawn by taking the distance of the particle from the source on X-axis and its displacement
on Y-axis.
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IIT - FOUNDATION - SET - V
Properties of the wave:
1) Wavelength () : The distance between two consecutive particles that are in the
same state of vibration or phase is known as the wavelength.
= A1A2 = B1B2 = C1C2. etc., .............
(in the above figure)
A portion of the wave of length equal to constitutes one wave.
2) Frequency (n) : The frequency of a wave represents the number of waves that
pass over a reference point in one second.
S.I Unit: hertz (or) cycles/sec.
The pitch of the sound depends on frequency. Low pitch sounds
are flat and high pitch sounds are shrill.
3) Time Period (T) : The time taken by the particles of the medium to complete one
oscillation. It is also the time taken by the source of sound to
complete one oscillation. T = 1/n.
4) Amplitude (A) : The maximum displacement of a particle of the effected medium
from its mean position is called the amplitude. The loudness of
sound is directly proportional to the square of the amplitude.
5) Velocity (V) : The relation between velocity (V), frequency (n) and
wavelength () of the wave is V = n
When a wave passes from one medium to another, then the frequency (n) remains
same but the velocity (V) and the wavelength () undergo a change.
Laplace’s formula for the velocity of sound (V) in a gas :
P RT
V= (or) V=
d M
where P is the pressure of the gas
d is the density of the gas
T is the absolute temperature of the gas
M is the molecular weight of the gas
44
PHYSICS
is the ratio of the molar specific heats of the gas
1: If V 1 and V 2 are t he velocit ies at temperatures T 1 and T 2 in a gas
V1 T1
respectively then, . V
2 T2
2: If V 1 and V 2 are t he velocities in the two gases of same and at same
V1 M2
temperature then, . V
2 M1
3: If the velocity in one gas at temperature T1 is the same as that in another gas at temperature T2 then,
T1 T
2 .
M1 M 2
Differences between Progressive and Stationary waves:
Progressive wave Stationary wave (Standing wave)
The progressive wave propagates in a Stationary waves are formed due to
medium without damping or superposition of two identical waves
obstruction. It can be transverse or moving in opposite directions.
longitudinal.
All the particles of the medium execute The particles of the medium execute
S.H.M with same amplitude and S.H.M with same frequency but with
frequency. different amplitudes.
When stationary waves are set up in
the medium then nodes and anti-nodes
are formed.
The equation of a progressive wave is The equation of a stationary wave is
y = A sin (t kx) y = 2 A coskx. sint
where y is the displacement ; where y is the displacement and
A is the amplitude. 2
2Acoskx = 2A cos x is the
2
2n and constant k
amplitude.
There is a flow of energy in the There is no net flow of energy in the
direction of wave propagation. medium.
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IIT - FOUNDATION - SET - V
1) Phase Difference : As a progressive wave is travelling in a medium the phase difference
between two particles that are at a distance ‘x’ apart is
2
= .x radian.
2) Nodes : Points where particles are at rest, in the stationary wave. The amplitude
2
2A cos x is zero at the nodes. So nodes appear at such positions
2
for which cos x0
3 5
i.e. nodes occur at the positions for which x , , ,......
4 4 4
3) Anti-nodes : Points where particles vibrate with maximum amplitude, in the propagation of a stationary
wave.
2
i.e. Antinodes occur at such positions for which the amplitude 2A cos x 2A
3
(or) x o , , , ,......
2 2
4) Distance between two consecutive nodes or anti-nodes = /2.
5) Distance between a node and an immediate anti-node = /4.
Free vibrations, Forced vibrations and Sympathetic vibrations:
1. Free vibrations: If a body vibrates with its natural frequency then its vibrations are called free
vibrations.
Ex 1: The vibrations produced in a tuning fork when it is vibrated and then left to itself.
1: The vibrations produced when a stretched wire is plucked and then left to itself.
2. Forced vibrations: When a body is subjected to external periodic force, then it vibrates with
the frequency of external periodic force. Such vibrations are called forced vibrations. Forced
vibrations are generally damped.
Ex 1: The vibrations of a microphone diaphragm due to speech.
2: The vibrations of a bridge, when a train passes over it.
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PHYSICS
3. Resonance (or) Sympathetic vibrations : When a body is made to vibrate under the action
of an external periodic force whose frequency is equal to the natural frequency of vibrating
body, then the body begins to vibrate with larger and larger amplitude. This is called resonance.
At resonance, energy transfer would be rapid, resulting in larger amplitude of the vibrations in
the body.
Ex : Two tuning forks having the same natural frequency are kept on two h o l l o w
wooden boxes as a shown in the diagram. If one is set into vibrations, then the other will be
found picking up vibrations of large amplitude. This will not be noticed, if the forks are of
different frequencies.
Velocity of transverse wave along a stretched string :
T
The velocity of a transverse wave in a stretched string is given by , V where ‘T’ is the
tension in the string and ‘ ’ is the mass per unit length of the string (or) linear density of the string.
1) Transverse stationary vibrations of strings : Let a string of length ‘ ’ be stretched between
the fixed ends A and B and plucked. It can have different modes of vibrations as shown below.
a) Fundamental mode (or) First harmonic: The string vibrates as a single loop with nodes at ‘A’ and
‘B’ as shown in figure.
2
2
But V = n V = n(2 )
T
We also know that V =
1 T
From and we get n where ‘n’ is the frequency of
2
fundamental mode.
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IIT - FOUNDATION - SET - V
b) First overtone (or) Second harmonic : In this mode the string vibrates with two loops as
shown in figure. Let and be the present wave length and frequency respectively.
Then 1
V 1 T T
n1 = V n 1 1 and V
1
1 T
n1 2
2
= 2n, where ‘’ is the frequency of the first overtone or second harmonic.
c) Second overtone (or) Third harmonic : In this mode, the string vibrates with three
loops as shown in figure. If the new wave length and frequency are then,
3 2
2 and n 2 then, ,
2
V V 3V
n2 = = =
2 2 / 3 2
1 T
n2 = 3
2
n 2 = 3n , where ‘’ is the frequency of the second overtone or third harmonic.
Thus, the frequencies n, ...........of different modes are in the ratio of natural numbers.
p T
If the string vibrates with ‘p’ loops then n
2
This is the frequency of p harmonic (or) (p-1)th overtone.
th
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PHYSICS
2) Laws of Vibrations of Stretched Strings :
The expression for the frequency (n) of the fundamental mode of vibration is given by n
1 T
= , where ‘l’ is the length of the vibrating segment, ‘T’ is the tension applied and ‘ ’ is the linear density
2
(or) mass per unit length of the string.
The following laws can be deduced from this expression:
For a given wire, held under constant tension, the frequency ‘n’ of the fundamental note is inversely
proportional to the length ‘ ’ of the vibrating segment.
1
n or n 1 1 = n 2 2 ; [ ‘T’ and ‘ ’ are constants. ]
For a given wire of given length, the frequency ‘n’ of its fundamental note is directly
proportional to the square root of the tension ‘T’ to which it is subjected.
n1 T1
n T or ; [ 'l' and ' ' are constants. ]
n2 T2
The frequency ‘n’ of the fundamental note emitted by a wire of constant length, subjected to a constant
tension is inversely proportional to the square root of its linear density or mass per unit length ‘ ’.
1 n
n or 1 2 ; [ ‘T’ and ‘l’ are constants. ]
n2 1
Experimental verification of the laws of transverse vibrations of the string using sonometer:
A) Description:
A sonometer consists of a wooden box (B) over which a stretched steel wire is arranged one end of this
wire (W) is tied to the nail (N) and the wire is passed over the movable knife edges (K1, K2) and a pulley (P).
From the other end a weight hanger with slotted weights of mass (M) is suspended. A paper rider (R) is placed
on the wire between k1 and k2 to feel the vibration of the wire.
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IIT - FOUNDATION - SET - V
B) Procedure:
Verification of Law - 1:
A tuning fork of frequency ‘n’ is excited and is placed vertically with its stem on the wooden box in
between k1 and k2. The distance between k1 and k2 is gradually increased without changing the tension (Mg) in
the wire. The paper rider flutters vigorously when the wire between k1 and k2 vibrates in resonance with the
tuning fork. At this stage, the length of the wire between k1 and k2 is measured as ‘l’. The frequency of this
segment of the wire is ‘n’ itself.
A number of observations are taken with tuning forks of different frequencies. The observations and
calculations are shown in the following tabular form.
S.No. Frequency of the Length of the vibrating n l = constant
tuning fork (n) segment(l)
The values in the last column are found to be equal to each other. This verifies the first law.
Verification of Law - 2 :
The frequency ‘n’ is to be kept constant for the verification of the second law. So the same
tuning fork is used throughout.
The tension ‘T’ in the wire is changed by increasing the hanging load ‘M’ in steps of 0.5kg and
for every tension the length of the wire that vibrates in resonance with the same tuning fork is measured. The
observations are entered into the following tabular form.
S.No. Load (M) Length of the vibrating l
segment(l) = constant
M
The values in the last column are found to be equal to each other. This verifies the second law
indirectly.
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PHYSICS
Verification of Law - 3 :
A number of wires of different thickness and made of different materials are selected. First, the
mass per unit length ( ) is determined for each wire. This is done by finding the mass of the wire with physical
balance and dividing the mass with its length. Keeping the tension constant, the resonating length (l) of each wire
is determined independently using the same tuning fork. The observations are entered into the following tabular
form.
S.No mass per unit length of Length of the vibrating l = constant
. the wire ( ) segment(l)
The values in the last column are found to be equal to each other. This verifies the third law
indirectly.
C) Precautions:
1) A rubber hammer only is to be used to excite the tuning fork.
2) The pulley must be frictionless.
3) The length (l) is measured only when the paper rider flutters with maximum amplitude.
1: Explanation for the verification of Law-2 (not for IPE)
The aim is to verify that n T . If the mass of the load suspended from the string
is M then T = Mg.
n1 n2
So the aim is to verify that n M (or) when l and are Constants.
M1 M2
But we may not find tuning forks of frequency n1 and n2 in the given box. So,
this law is verified indirectly. Using a single fork of frequency ‘n’ we shall change
the tension from T1 = M1g to T2 = M2g and find out the resonating lengths of the
same wire as ‘l1’ and ‘l2’.
We shall now apply law -1 and find out the frequencies and for any
given length ‘l’ in these two cases.
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IIT - FOUNDATION - SET - V
Case 1 : Tension T1 = M1g, length = l1, frequency = n, expected length = l,
expected frequency = n 1 .
nl 1
From law-1, it follows that nl1 = l (or) n1
Case 2 : Similarly, for tension T2 = M2g, the expected frequency for the same
Length ‘l’ of same wire is given by
nl 2
n2
nl1 nl 2
n1 n2
l l
So verification of is equal to = =
M1 M2 M1 M2
l1 l2 l
(or) or
M1 M2 M = constant.
2 : Explanation for the verification of Law-3 (not for IPE)
1
The aim is to verify that n (or) n 1 1 n 2 2 when T and l are constants.
But we may not have tuning forks of frequency and with us. So this law is verified indirectly using a single
fork of frequency ‘n’.
We shall first select two wires of different linear densities and and subject them to same tension. Using a
fork of frequency ‘n’, the resonating lengths of these wires are determined as ‘l1’ and ‘l2’.
Applying law-1, we shall now calculate the frequencies and for the same length ‘l’ in these two cases.
Case : Linear density = , length = l1, frequency = n, expected length = l,
expected frequency = n
1
nl1
from law-1, nl= l. (or) n 1
l
Case : Linear density= 2 , length = l2, frequency = n, expected length = l,
expected frequency = n1
nl 2
from law-1, n2
l
n l1 nl
So to verify n 1 1 n 2 2 is equal to 1 = 2 2
l l
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PHYSICS
(or) l 1 1 l 2 2
(or) l constant
1: Production:
When two sources of nearly the same frequency are sounded together then the resulting sound waxes and
wanes. This waxing and waning in the loudness of the sound is called beats.
The number of times it waxes in one second (or) wanes in one second is called the
Beat frequency. If n1 and n2 are the frequencies of the sources, then beat frequency is equal to ( ~ ). The
persistence of hearing is 0.1 second. So beats can be heard only when the interval between the successive beats
is not less than 0.1 second. In other words, beats can be heard only when the beat frequency is less than 10 or
( ~ ) is less than 10.
2: Explanation:
The sound waves from the two sources have same amplitude but have different frequencies ‘’ and .
These waves mixup and form a resultant wave. As a result of this interference, the amplitude of the resultant wave
changes continuously with time with alternate minimum and maximum amplitudes.
Minimum amplitude = a - a = 0
Maximum amplitude = a + a = 2a
Correspondingly minimum loudness = 0
maximum loudness = 4a 2
The resultant wave is as shown below.
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IIT - FOUNDATION - SET - V
1
Time interval between successive maxima or minima is n n seconds.
1 2
Beat frequency = ( n 1 n 2 )
3: Uses (or) Applications of Beats:
1) To standardise the frequency of a given tuning fork: The prongs of the given tuning fork
are slowly filed so that the beats subside while vibrating along with a standard tuning
fork.
2) To tune the Musical instruments : Musical instruments like veena etc. can be tuned
to a standard tuning fork by eliminating beats. The tension in the wire is adjusted until
the beats subside while vibrating along with standard fork.
3) To detect the release of dangerous gases while digging in the mines: A column of air in
a resonator is set to vibrate in resonance with a tuning fork. Any contamination of the air
changes the frequency of the resonator slightly. So beats are produced. This amounts to
raising the alarm.
4: Theory of Beats : (not for IPE)
Let two sources of sound of nearly the same frequency ‘’ and ‘’ be excited. The displacements ‘y1’ and
‘y2’ of a particle due to the waves from these two sources are
y1 = asin2n1t and y2 = asin2n2t
The net displacement y of this particle is given by
y = y1 + y2
y = a sin2t + asin2t
2 n1 n2 t n n
y = 2a.cos .sin 2 1 2 t
2 2
n1 n2
So the resultant sound has a frequency of
2
n1 n2
The amplitude of the resultant wave is given by A = 2a cos 2 t
2
n1 n2
So the amplitude is changing simple harmonically with a frequency
2
The loudness A2
So the loudness of the sound changes with time
Loudness = 4a 2 .cos 2 [ ( n1 n2 ) t ]
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PHYSICS
(a) Waxing with maximum loudness of occur when
1 2 3
t=0 , , .........
n1 n 2 n1 n 2 n1 n 2
(b) Waning with minimum loudness equal to zero occur when
1 3 5
t = 2 n n , 2 n n , 2 n n ..........
1 2 1 2 1 2
Time interval between the successive maxima or minima in given by
1
T=
(n1 n2 )
Beat frequency = n1 n 2
Doppler Effect:
Statement : The apparent change in the frequency of a sound wave, due to the relative motion
between the source and the observer is called the Doppler Effect.
Derivation :
Step 1: Change in wavelength of the wave:
The sound wave from a stationary source of frequency ‘n’ travels with a velocity ‘V’ in still air. The
V
wavelength of this wave is given by =
n
The wavelength of the wave is modified due to both the blowing of wind and the motion of the source.
(a) Effect of wind :
If the wind is blowing with a velocity ‘W’ then the velocity of the wave in the direction of wind is (V+W)
and the velocity in the opposite direction is (V-W)
So in the direction of the wind, the ‘n’ waves that are produced would occupy a distance (V+W).
(b) Effect of the motion of the source :
Let the source be moving in the direction of the wind with a velocity . The velocity of the wave in the
direction of the wind is still (V+W) as it is not effected by the motion of the source. But in the direction of the wind
the n-waves that are produced in
1 second, would occupy a distance (V+W-).
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IIT - FOUNDATION - SET - V
The wavelength of the wave travelling in the direction of the wind is given by
V W VS
1
n
Step 2 : Pitch or frequency of the sound felt by the observer :
The frequency of the sound, felt by the observer depends on
(a) the wavelength of the wave approaching the observer and
(b) the velocity with which the wave approaches him.
The wave travelling in the direction of the wind approaches the observer with a
velocity (V+W). But if the observer is also moving in the same direction with velocity then the
wave approaches him with a velocity (V+W-)
frequency of the sound Velocity of approach
=
as felt by observer Wavelength
(V W V0 )
n1
V W VS
n
VWV0
n1 n
Apparent frequency = V WV
S
Case 1 : When there is no blowing of wind then W 0
V V0
n1 n
V VS
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PHYSICS
Case 2 : When the observer is at rest and the source is moving towards the observer
1 V
then n n
V VS
Case 3 : When the source is at rest and the observer is moving away from the source -
V V0
then n1 n
V
Case 4 : When the source is moving in the same direction and the observer is moving in
opposite direction i.e. the source and the observer are moving towards each
other.
V V0
then n1 n
V VS
Case 5 : When the source and the observer move away from each other
V V0
then n1 n
V VS
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IIT - FOUNDATION - SET - V
Sign convention:
a) Wave velocity ‘V’ is always taken as positive.
b) Any other velocity is taken as positive if it is along the direction of ‘V’ and is taken negative
if it is in the opposite direction.
Limitations of Doppler effect:
1) Doppler effect could be felt only when the velocity of the observer is less than the velocity of
the wave that is approaching him.
2) Also Doppler effect cannot be felt at the instant when the source is moving relative to the
observer in a direction at right angles to the direction in which he is receiving the sound wave.
Applications of Doppler effect :
1) Doppler effect is used by the traffic police for estimating the speed of a running vehicle.
2) This principle is used in tracking the Earth’s satellite.
3) This effect is used in accurate navigation of aircraft and it is the basis of accurate target
bombing techniques developed recently.
4) Motion of the Star : The spectrum of light from a particular star can be formed. the shift in
the spectral lines towards red end of the spectrum indicates an increase in the wavelength and
decrease in the frequency of the light. This shows that this star is moving away from the Earth.
Similarly, the shift in the spectral lines towards violet end indicates that the star is moving
towards the Earth.
5) Saturn rings : Observing the spectra of light from the inner and outer edges of the Saturn
rings, it is concluded that the rings constitute swarm of particles and not hard discs.
6) Double Stars : The spectral lines of the light from certain stars become alternately ‘single’
and ‘double’. This is because the star under observation actually consists of two stars which
are rotating about common axis. When both the stars are along the line of sight only one
spectral line is seen. But when one star is seen approaching and the other receding, two lines
are seen.
Acoustics:
1. Echo: The sound of the wave reflected from a distant reflecting surface is called the echo.
The persistence of hearing for a human ear is 0.1sec. So the echo of a sharp sound like
clapping can be felt only when it arrives after 0.1sec. If the velocity of sound in air is
assumed to be 340m/s then the distance travelled by the waves in 0.1second would be 34m.
58
PHYSICS
Considering the to and fro journey of the wave, the reflecting surface should be at a minimum
distance of 17m from the person, so that he could hear the echo clearly.
Applications of echoes : The echo of the ultrasonic sounds is used in (1) estimating the
height and position of the aeroplane and (2) in finding the depth of the oceans (SONAR).
2. Reverberation : When a sharp sound is produced in a large hall or auditorium then a person
in the hall continues to hear it for a long period due to multiple reflections from the surrounding
walls, ceiling and floor of the hall. The sound energy gets distributed throughout the hall due to
multiple reflections. This kind of distribution of sound energy is called reverberation.
3. Time of Reverberation : The one effect of reverberation, specially in large halls is that the
loudness of a sharp sound slowly increases to a maximum due to multiple reflections from
distant walls and then gradually decreases due to absorption. In other words, the time of
reverberation becomes too large, specially for large halls.
Definition of time of reverberation : The time for which the sound continues to be
heard in the room, even after the cutting off the source is called the time of reverberation of the
room.
Standard definition of the time of reverberation: The time of reverberation of a hall
is the time taken, after cutting off the source, for the sound to fall in intensity to 10-6 of its initial
value.
4. Sabine’s formula: Sabine gave the following empirical formula for the time of reverberation
after a long experimental study.
kV
Time of reverberation, T = Where V = volume of the hall,
as
s = area of the reflecting surface and
a = absorption coefficient of the surface.
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IIT - FOUNDATION - SET - V
k .V
If there are a number of different reflecting surfaces, then T=
as
The constant k is 0.171 in S.I. system.
The absorption coefficient of a surface is defined as the ratio of the energy absorbed by it to the energy
incident on it. An open window reflects no sound. So an open window is assumed to be a perfect absorber and
the absorption coefficient is redefined.
Definition of Absorption Coefficient: The absorption coefficient of a surface is defined as the ratio of
the sound energy absorbed by it to that absorbed by an equal area of an open window.
(1) Experimental determination of the absorption coefficient by stationary wave method:
Sound waves of amplitude from the source are made to incident on the test materials.
The reflected waves have an amplitude which is less than due to absorption. The incident and
reflected waves get superposed and produce stationary waves.
Amplitude at the node (N)= a a
1 2
Amplitude at the anti-node [AN)= a1 a2
A microphone is moved from the source to the test material and the current from it is measured. The
maximum and minimum values of these currents are noted as i1 and i2.
i1 a 1 a 2
i 2 a1 a 2
a1 i1 i2
(1)
a2 i1 i2
2
incident energy = a1
2
reflected energy = a2
2 2
absorbed energy = a1 a2
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PHYSICS
a12 a22 a2 2
Absorption coefficient = 1 (2)
a12 a12
From equations (1) and (2) absorption coefficient =
Thus, absorption coefficient can be determined by knowing i1 and i2.
(2) Experimental determination of the absorption coefficient by Reverberation method:
Special enclosures called Reverberation chambers are constructed for this method. The surfaces of this
chamber reflect more sound. Without any absorbing materials in the chamber, the reverberation time T1 is
determined. Then the material whose absorption coefficient is to be determined is pasted on the walls of the
chamber and the time of reverberation T2 is determined. Then absorption coefficient of the material is determined
by
0.17V 1 1
a= ao
S T2 T1
Where V = volume of the chamber
S = surface area of the walls
ao= absorption coefficient of the chamber before pasting.
Theory: By Sabine’s formula,
1 a0 S . 1
a) Initially, T 0.17 V ………………(1)
1
1 aS . 1
b) After pasting, T 0.17 V …………..(2)
2
1 1 aS . 1 a S 1
0 .
T2 T1 0.17 V 0.17 V
S .1
= (a a 0 )
0.17 V
0.17 V 1 1
a a0
S T2 T1
0.17 V 1 1
a= a 0
S T2 T1
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IIT - FOUNDATION - SET - V
Acoustic Design or the main conditions that are to be fulfilled for the better acoustic quality of a building:
One has to take utmost care in the design of an auditorium or a concert hall so that
(a) any speech delivered in the hall is intelligible to all the audience and
(b) the music which is played in the hall should not loose its charm.
The following are the important features that determine the acoustic quality of the building.
1) Optimum Reverberation: The one good effect of reverberation is that it helps the person
even in the back row to receive the sound in sufficient loudness. In the absence of reverberation,
the appeal will be very poor and the room appears dead. This is the case with the open-air
theatre.
At the same time if the time of reverberation of the room is large, there will be overlapping
of the syllables in the hearing of speech. So the speech becomes unintelligible. In the worst
case, the speech is lost entirely.
For the better acoustics, the optimum value of the time of reverberation is found to be 1
to 2 seconds for music and 0.5 to 1 seconds for speech. From the Sabine’s law it follows that
the time of reverberation becomes very large for large auditorium. However, it can be reduced
to optimum by properly covering the walls and ceiling with pulp boards which absorb the
sound and suppress multiple reflections.
2) Focussing effect: Walls or roof with spherical and cylindrical curvatures cause bad focussing.
Parabolic curvature for the wall behind the stage with the speaker at the focus is desirable.
This distributes the sound uniformly throughout the auditorium.
3) Extraneous sounds: Entry of sound into the auditorium, from the external sources should be
cut off.
4) Echelon effect: Reflection of sound from a flight of stairs produces a kind of musical sound
that disturbs the original sound. In such a case, the stairs should be covered with a carpet.
62
PHYSICS
ASSIGNMENT
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IIT - FOUNDATION - SET - V
1. The velocity of sound in air and in water are 350 m/s. and 1470 m/s. If the wave length in air is 1m, the
frequency and wave length in water are
1) 350 Hz, 4.2m 2) 700 Hz, 4.2m
3) 1470 Hz, 4.2m 4) 1470 Hz, 1m
2. The vibrations of a string of length 60 cm fixed at both ends are represented by the equation y = 4 sin
x
cos (96 pt) where x and y are in centimetres and ‘t’ is in seconds. The maximum displacement
15
of a point at x=5 cm is
1) 2 3 cm 2) 2 cm 3) 4 3 cm 4) None
3. The vibrations of a string of length 60cm fixed at both ends are represented by the equation, y = 4sin
x
cos (96 pt), where x and y are in cm and t in seconds. The nodes are located along the string at
15
1) 0, 15, 30, 45, 60 2) 0, 20, 40, 60
3) 0, 10, 20, 30, 40, 50, 60 4) 0, 20, 40, 60
x
4. A simple harmonic wave is represented by the relation, y (x, t) = a0 sin2p (ft - ). If the maximum
particle velocity is three times the wave velocity, the wave length l of the wave is
1) a 0 3 2) 2a 0 3 3) pa0 4) a 0 2
t x
5. The equation of a progressive wave is y = 0.5 sin 2 where x is in cm. The phase
0.004 50
difference between two points separated by a distance 10 cm at any instant is
1) 720 2) 360 3) 900 4) 480
6. A sonometer wire of length 100 cm is divided into three segments, whose fundamental frequencies are
in the ratio 20 : 5 : 4. The distance between the two bridges is
1) 20 cm 2) 30 cm 3) 60 cm 4) 40 cm
7. The linear density of a vibrating string is 1.3 ´ 10-4 kg/m. A transverse wave is propagating on the string
and is described by equation y = 0.021 sin(x+30t) where x and y are in meters and t in seconds. The
tension in the string is nearly
1) 0.12 N 2) 0.48 N 3) 1.2 N 4) 4.8 N
8. An addition of 24kg to the tension of a string, changes the frequency of the string to three times the
original frequency. The original tension is
1) 2kg wt 2) 4kg wt 3) 6 kg wt 4) 3kg wt
9. The length of a string is reduced by 25% and tension is increased by 125%. Its fundamental frequency
increases by
1) 50% 2) 100% 3) 150% 4) 200%
64
PHYSICS
10. Two strings A and B made of the same material have equal length, the cross sectional area of A is half
that of B, while the tension on A is twice that on B. The ratio of the velocities of transverse waves in A
and B is
1) 2 : 1 2) 2 : 1 3) 1 : 2 4) 1 : 2
11. A metal wire is held at the two ends by a rigid support. At 900C the wire is just taut with negligible
tension. If Y = 1.8 x 1011N/m2, = 18 x 10-6/0C, and d = 9 x 103 kg/m3 the speed of the transverse wave
in this wire at 00C is
1) 20m/s 2) 45m/s 3) 90m/s 4) 180m/s
12. The stress developed in a stretched wire of length 1m is 20 x 108Pa. If the density of the metal of the
wire is 8000 kg/m3 then the frequency of its second overtone is
1) 750Hz 2) 500Hz 3) 250Hz 4) 1500 Hz
13. A string of linear density 0.03 g/cm vibrates in three loops to emit a frequency of 600Hz while it is under
a tension of 480 N. Length of the string is
1) 2 m 2) 3 m 3) 0.5 m 4) 1 m
14. A wire having a linear density 0.1 kg/m is kept under a tension 490N. It is observed that it resonates at
a frequency of 400Hz and the next higher frequency 450Hz. The length of the wire is
1) 0.4m 2) 0.7m 3) 0.6m 4) 0.49m
15. A stretched string of one meter length and mass 5 ´ 10-4 kg is fixed at both ends and under a tension of
20N. When plucked at a distance of 25cm from one end, the string would vibrate with a frequency of
1) 400 Hz 2) 100 Hz 3) 200 Hz 4) 256 Hz
16. The length, diameter, tension and density of string B are double to that of String A. The overtone of B
that is in unison with 1st harmonic of A is
1) 1st 2) 2nd 3) 3rd 4) 4th
17. A uniform rope of mass 0.1kg and length 2.45m hangs from a ceiling. Then the speed of the transverse
wave in the rope at a point 0.5m distant from the lower end is (g=10ms-2)
1) 2 m/s 2) 5 m/s 3) 10 ms-1 4) 7 ms-1
18. A wire of density 5gcm-3 is stretched between two clamps 100cm apart while subjected to an exten-
sion of 0.5cm. If young’s modulus of the wire is 9x1011 dyne cm-2 then the lowest frequency of trans-
verse vibrations in the wire is
1) 100Hz 2) 150Hz 3) 200Hz 4) 300Hz
19. A tuning fork ‘x’ gives 4 beats per second with a tuning fork of frequency 256 Hz. When ‘x’ is loaded
with wax, it again gives 4 beats/s. The frequency of x is
1) 256 Hz 2) 252 Hz 3) 262 Hz 4) 260 Hz
20. A set of 25 tuning forks are arranged in a series of decreasing frequencies. Each fork gives 3 beats per
second with the succeeding one. The first fork is octave of the last. The frequency of the sixteenth fork
in series is
1) 77Hz 2) 99Hz 3) 144Hz 4) 155Hz
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IIT - FOUNDATION - SET - V
21. A disc siren has 20 holes around its rim. When revolving at 1020 r.p.m it gives 6 beats per second with
a standard fork. If now the speed of revolution of the siren is increased to 1050 r.p.m number of beats
heard with the same fork is 4 beats/sec. Frequency of standard fork is
1) 340 Hz 2) 346 Hz 3) 350 Hz 4) 334 Hz
22. Two identical stringed instruments have a frequency of 100Hz. If tension in one of them is increased by
4% and they are sounded together then number of beats produced in one second is
1) 1 2) 8 3) 4 4) 2
23. Two tuning forks A and B produce 6 beats per second. A resonates with a string of length 64.6 cm and
B resonates with a string of length 64 cm. Then frequencies of A and B are
1) 640, 646 Hz 2) 620, 626 Hz 3) 650, 656 Hz 4) 630, 636 Hz
24. A knife edge separated the sonometer wire of length 100cm, into two parts, which differ in length by
2cm. If the two parts produce 6 beats per second, while vibrating together the frequency of the smaller
portion is
1) 306Hz 2) 153Hz 3) 147Hz 4) 294Hz
25. A sonometer wire produces 4 beats per second with a tuning fork when the length is either 204cm or
208cm. Then the frequency of the tuning fork is[2]
1) 400Hz 2) 412Hz 3) 388Hz 4) 206Hz
26. A tuning fork of frequency 480Hz produces 10 beats / second when sounded with a vibrating string. If
a slight increase in tension produces fewer beats per second than before, then the original frequency of
the string was.
1) 460 Hz 2) 480Hz 3) 490 Hz 4) 470 Hz
27. Three forks A, B and C have frequencies such that the frequency of A is 0.7% greater than that of B and
the frequency of C is 0.7% less than that of B. If A and C produce 56 beats in 10 seconds, the
frequencies of the forks A, B and C are.
1) 428 Hz, 400Hz, 372Hz 2) 402.8Hz, 400Hz, 397.2Hz
3) 40.3Hz, 40Hz, 39.7Hz 4) 397.2Hz, 408.8Hz, 406Hz
28. A stretched sonometer wire is in unison with a tuning fork. When the length of the wire is increased by
2%, the number of beats produced per second is 6. Frequency of the tuning fork is
1) 300Hz 2) 306Hz 3) 600Hz 4) 606Hz
29. A sonometer wire vibrates with a frequency 1200Hz. When the body suspended from the wire is kept
immersed in water, the frequency changes to 800Hz. The specific gravity of the material of the body is
1) 1.5 2) 2.4 3) 1.8 4) 2.6
30. Two progressive waves Y1 = 5 sin 400 pt and Y2 = 4 sin 404 pt moving in the same direction superpose
on each other producing beats. Then the number of beats per second and the ratio of maximum and
minimum intensity of the resultant waves are respectively
1) 2 and 9 : 1 2) 2 and 81 : 1 3) 4 and 9 : 1 4) 4 and 81 : 1
31. A source of sound of frequency 500Hz is moving towards north with a velocity 10m/s. Wind is blowing
towards south with a velocity 5 m/s. Velocity of sound in still air is 340m/s. The wavelength of the wave
travelling towards North is
1) 0.65m 2) 0.71m 3) 0.69m 4) 0.67m
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PHYSICS
32. A train moves towards a stationary observer with speed 34m/sec. The train sounds a whistle and its
frequency felt by the observer is n1. If the train speed is reduced to 17m/sec, the frequency felt by him
is n2. If the speed of sound is 340m/s, then n1 : n2 is
1)18 :19 2) 7 :6 3) 6 :7 4) 19 : 18
33. Two express trains travelling at 72 kmph approach each other whistling. If the frequency of the note
emitted is 800 Hz, the apparent frequency of sound as heard by a passenger in the other train is (Veloc-
ity of sound in air is 340 m/s)
1) 920 Hz 2) 900 Hz 3) 820 Hz 4) 800 Hz
34. A person is watching two trains, one leaving and the other coming in, with equal speed of 4 m/s. If they
sound their whistles, each of natural frequency 240 Hz, the number of beats per second heard by the
person will be nearly equal to
(velocity of sound in air =320m/s)
1) 3 2) 2 3) zero 4) 6
35. Two sources A and B, which are at certain distance apart produce sounds of frequency 680Hz. A
person, while walking from A to B hears 5 beats per second. If the velocity of sound in air is 340m/s
then the velocity of the person is
1) 2.5m/s 2) 2m/s 3) 3.5m/s 4) 1.25m/s
36. Two sources A and B produce sounds of same frequency. A person running from A to B hears 4beats/
s. If the frequency of each source increases by 100Hz then 6 beats/s are heard. The original frequency
of the sources is
1) 100 Hz 2) 200 Hz 3) 300 Hz 4) 400 Hz
37. When a train is crossing an observer the ratio of the apparent frequencies of the whistle of the train heard
by him is 6:5. The velocity of sound in air is 352 m/s. Then the velocity of the train is
1) 25m/s 2) 35m/s 3) 32 m/s 4) 30m/s
38. A car with a horn of frequency 620Hz, travels towards a large wall at a speed of 20m/s. If the velocity
of sound is 330 m/s, the frequency of the echo of the sound horn as heard by the driver of the car is.
1) 700 Hz 2) 657.57Hz 3) 549.1 Hz 4) 175 Hz
39. A whistle emitting a sound of frequency 600Hz is tied to a string of length 1.5m and rotated with an
angular velocity 20 rad/s in the horizontal plane. Then the range of frequencies heard by an observer
stationed at a large distance from the whistle is (velocity of sound in air= 330m/s)
1) 500Hz to 650Hz 2) 550Hz to 650Hz 3) 550 Hz to 660Hz 4) 500Hz to 660Hz.
40. A girl swings in a cradle with period p/4 seconds and amplitude 2m. A boy standing infront of it blows
a whistle of natural frequency 1000 Hz. The minimum frequency as heard by the girl is (velocity of
sound in air = 320 ms-1)
1) 850 Hz 2) 1000 Hz 3) 750 Hz 4) 950 Hz
41. A body emitting sound of frequency 350Hz is dropped from a balloon raising vertically upwards with
constant velocity 5 m/s. Frequency of sound as felt by the observer in the balloon 2s after the release is
: (Velocity of sound in air is 335 m/s & acceleration due to gravity is 10 ms-2 )
1) 330Hz 2) 335Hz 3) 340Hz 4) 355Hz
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42. A whistling railway engine approaches a statioanry observer on the platform with speed 34 m/s and
recedes away with same speed. If the difference in the two apparent freqnecies heard by the observer
is 200 Hz and velocity of sound in air is 340 m/s, the real frequency of the sound emitted by the whistle
is
1) 990 Hz 2) 1980 Hz 3) 495 Hz 4) 660 Hz
43. The apparent frequency as heard by the observer is 25% more than the real frequency of the source and
‘C’ is the velocity of sound in air. Then
1) source approaches stationary observer with speed C/5 only
2) observer approaches stationary source with speed C/4 only
3) source approaches with speed 2C/5 but observer recedes away with speed C/4 only
4) any of the above
44. A man is driving at 72 km/hr on a straight road heading towards a hill. He sounds the horn and hears its
echo 4 seconds afterwards. Assume the speed of sound in air to be 340 m/s. The distance from the hill
at which the horn was sounded is
1) 320 m 2) 420 m 3) 720 m 4) None
45. A man stationed between two cliff’s fires a gun. He hears the first echo after 2 seconds and the next
after 5 seconds. If the velocity of sound in air is 340 m/s, the distance between the two cliffs is
1) 1.19 km 2) 2.19 km 3) 1 km 4) 2 km
46. In a room where sound intensity is not uniform absorption of sound by a material depends on
1) the shape of the material 2) size of it’s material
3) position of the material 4) all of the above
47. An auditorium has dimension 100 x 10 x 4 m. The absorption coefficient of the wall and the roof is 0.2.
Open windows cover an area 20 m2. Absorption coefficient of the floor is 0.164. The time of rever-
beration is
1) 1.0 s 2) 0.97 s 3) 1.02 s 4) 1.1 s
48. A room of dimensions (100 x 50 x 10)m3 has the reverberation time 5 seconds. After finding the
absorption of the room, if 680 persons occupy the room each having sound absorption of 0.5 metric
sabine, then the new reverberation time approximately...
1) 2.15s 2) 3.15s 3) 4.167s 4) 5.15s
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PHYSICS
KEY
1 1 16 3 31 1 46 4
2 1 17 2 32 4 47 2
3 1 18 2 33 2 48 3
4 2 19 4 34 4
5 1 20 2 35 4
6 4 21 2 36 2
7 1 22 4 37 3
8 4 23 1 38 1
9 2 24 3 39 3
10 1 25 2 40 4
11 4 26 4 41 1
12 1 27 2 42 1
13 4 28 2 43 4
14 2 29 3 44 3
15 3 30 2 45 1
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4 MAGNETISM
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PHYSICS
COULOMB’S INVERSE SQUARE LAW IN MAGNETISM:
Statement: The force of attraction or repulsion between two magnetic poles is directly proportional
to the product of their pole strengths and inversely proportional to the square of the distance between them.
Explanation: If two magnetic poles of strengths m1 and m2 are separated by a distance‘d’ in
vacuum or in air then the force (F) between them is given by
m1 m2 0 m1m 2
F F=
d 2
4 d 2
Where 0 is the permeability of the free space or air & 0 = 4 10-7 Hm-1 (Hm-1 stands for henry
per metre).
m1 m2
If the poles are situated in any medium other than vacuum then F = .
4 d2
0 r m1m 2
(or) F =
4 d 2
Where is the absolute permeability of the medium and r is the relative permeability of the
medium.
Relative permeability: It is the ratio between the absolute permeability of the medium and the
permeability of the free space. r
0 .
S.I. unit of pole strength is ampere - metre (Am).
Definition of ampere-metre: It is the strength of that pole which when placed in vacuum at a distance
of one metre from an identical pole repels it with a force of 10-7 newton.
Note : The pole strength ‘m’ is taken positive for N-pole and negative for S - pole.
Magnetic Field: It is the space around a magnet where its influence is felt.
Magnetic induction (induced field intensity) (B): The magnetic induction at a point is defined as the
force experienced by a unit north pole kept at that point.
The S.I. unit of magnetic induction is Tesla.
Definition of Tesla : The magnetic induction at a point is said to be one ‘’tesla’’, if a N-pole of
one ‘Am’ placed at that point experiences a force of one ‘newton’.
1) Magnetic induction field ‘B’, is a vector quantity.
2) For historical reasons the unit tesla is also frequently called either as
newton per ampere metre or weber per metre square.
1 tesla = 1NA-1 m-1 = wb m-2.
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3) A pole of strength ‘m’ situated in a magnetic induction field ‘B’ experiences a force F = mB or
F mB .
Moment of Couple acting on a magnet placed in a uniform magnetic field:
NS is a bar magnet of pole strength ‘m ’ and of length 2l. This magnet is placed in a
uniform magnetic field of induction ‘B’ with its axis making an angle with the field. A force
‘mB’ acts on N-pole along the field. A force of ‘mB’ acts on S-pole in the opposite direction.
These two forces are anti parallel, equal and non-collinear and hence constitute a couple.
Moment of couple C = Force perpendicular distance between the forces
C = mB SP
SP SP
= mB 2l sin In the le NPS , sin
SN 2l
C = (m 2l) B sin
The product of pole strength ‘m’ and length ‘2l’ of the magnet is called the magnetic moment and is
represented by M.
C = M B sin
If B = 1 T and = 900 then M = C.
Definition of Magnetic moment (M): It is the couple experienced by the magnet when it is placed in a
uniform magnetic field of one tesla with its axis at right angles to the direction of the magnetic field.
Note : Magnetic moment is a vector quantity.
The direction of the magnetic moment is from South Pole of the magnet to its north pole.
S.I. unit of magnetic moment is ampre-metre2 (Am2)
SHM of a magnet suspended in uniform magnetic field:
When a magnet of magnetic moment M is suspended in a uniform magnetic field of induction B, it comes
to rest in the direction of the field. This rest position of the magnet is called its mean position.
If this magnet is given a small angular displacement radian, then it experience a couple given by
C = MB Sin.
When is very small sin =
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PHYSICS
So, couple C = MB
(1)
Under the action of this couple as shown in the fig. the magnet tends to deflect back towards the mean
position with an angular acceleration ‘’ given by
C = I (2), where I is the moment of inertia of the magnet about the axis about
which it is oscillating.
From equation (1) and (2), I = MB
= ; is constant for a given magnet.
So, (angular acceleration) (angular displacement)
Hence, the magnet executes simple harmonic motion.
2 I
T 2
The period of its oscillation is given by MB MB
I
Magnetic induction ‘B’ due to a bar magnet:
a) The magnetic induction ‘B’ at a point ‘P’ on the axis of the magnet:
If M = magnetic moment
l = half the length of the magnet
d = distance of the point ‘P’ from the midpoint ‘O’ of the magnet
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IIT - FOUNDATION - SET - V
B = magnetic induction at ‘P’.
0 2 Md
Then, B = 4 d 2 l 2 2 tesla, in the S-N direction.
0 2 M
For a short bar magnet, d>>l, so B =
4 d 3
b)The magnetic induction at a point on the equatorial line or neutral axis of a bar magnet: The
magnetic induction at a point ‘P’ on the equatorial line is given by
0 M
B = 4 d 2 l 2 3/2 tesla, along N-S direction
0 M
For a short bar magnet d>>l, so B =
4 d 3
Magnetic lines of force:
A magnetic line of force is a curve along which a north pole would travel if it were free to move.
Properties of Magnetic lines of force:
Magnetic lines of force are closed curves.
A magnetic line of force starts from the North pole and ends with the South pole outside the
magnet, and within the magnet it is from South pole to North pole.
No two magnetic lines of force intersect each other.
The magnetic lines of force are closely spaced at the region where the intensity is high.
The tangent drawn at any point to a magnetic line of force gives the direction of the magnetic field
at that point.
Magnetic induction or Magnetic flux density (B):
The number of lines of force passing through unit area of a surface, held at right angles to the magnetic
74
PHYSICS
field is known as magnetic flux density B.
Magnetic flux through an area “A” is given by
= B . A = BA cos,
where = angle between B and the normal to the surface. Unit of magnetic flux: weber.
If = 00 then = B.A (or) B = Weber/metre2
A
= Magnetic flux through unit area
= Magnetic flux density
So the magnetic induction is called as the magnetic flux density.
Superposition of Magnetic fields: (not for IPE)
a) Mapping of the combined magnetic field of the earth and a bar magnet with North pole
facing Geographic North: A short bar magnet is placed horizontally in the magnetic meridian with
its North pole pointing Geographic North, on a white paper fixed to a drawing board. The magnetic
lines are drawn using a compass. Then, two null points are obtained on the equatorial line of the bar-
magnet as shown in the figure. Let d be the distance of each null point from the centre of the magnet
and BH be the horizontal component of the earth’s magnetic induction. The component BH is cancelled
by the induction B of the magnet at the neutral point, then
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IIT - FOUNDATION - SET - V
where M = magnetic moment of the bar magnet.
b) Mapping of the combined magnetic field of the Earth and a bar magnet with south pole facing
Geographic North: A short bar magnet is placed horizontally in the magnetic meridian with its
South pole pointing Geographic North, on a white paper fixed to a drawing board. The magnetic
lines are drawn using a compass. Then, two null points are obtained on the axial line of the bar-
magnet as shown in the figure. If d is the distance of each null point from the centre of magnet and BH
is the horizontal component of the earth’s magnetic induction, then BH = (where M = Magnetic
moment of the bar magnet)
Deflection magnetometer - comparison of magnetic moments:
Principle:
Tangent Law: When a magnetic needle is freely suspended in a region where there are two uniform
magnetic fields ‘B’ and ‘BH’ acting at right angles, then needle comes to rest with the axis making an
angle with BH so that B = BH tan q (BH is the horizontal component of earth’s magnetic induction
and B is the field produced by a magnet).
Explanation: The magnet comes to rest with its axis along the resultant of B and BH such that
B
Tan =
BH or B = BH tan
BH
N
S B
Method 1 :
Comparison of Magnetic moments of two bar magnets using a Deflection Magnetometer in tan-
A position by equal distance method:
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PHYSICS
a) Principle: Deflection magnetometer makes use of tangent law.
“ When a magnetic needle is freely suspended in a region where there are two uniform magnetic
fields ‘B’ and ‘BH’ acting at right angles, then needle comes to rest with the axis making an angle
with BH so that B = BH tan ”
b) Description : The deflection magnetometer consists of asmall magnetic needle, pivoted on a
pointed support. A light, thin and long aluminium pointer is attached at right angles to the
magnetic needle. The pointer moves over a scale graduated in degrees as 0 -90 - 00 -900 - 00.
0 0
A plane mirror is placed below the scale to observe the deflections without parallax error. The
magnetic needle, pointer and scale with mirror are enclosed in a box with a glass top called compass
box. It is kept at the central portion of a wooden bench over which a scale is attached.
c) Arrangement of the deflection magnetometer in tan-A position: The wooden bench is
rotated such that the aluminium pointer becomes parallel to it. i.e., the wooden bench is set in East-
West direction. Without disturbing the bench, the magnetometer is rotated until the aluminium pointer
reads 00-00. This is called tan-A position.
d) Procedure: One of the magnets of moment M1 is placed at a distance ‘d’ from the centre of
deflection magnetometer on one of the arms with its length parallel to the bench. The deflections are
noted as 1 and 2.The magnet is reversed pole to pole in its position and the deflections are noted
as 3 and 4. Four more such deflections are noted as 5, 6, 7 and 8 keeping the magnet on the
other arm and at the same distance. If the average of these deflections is 1then
0 2M 1
. B H Tan1
(1)
4 d 3
Similarly the average deflection for a second magnet of moment M2 and placed at the same
distance d is obtained as 2.
Then, 0 . 2M 2 B Tan (2)
3 H 2
4 d
M 1 Tan1
From equations and
M 2 Tan 2
A number of observations can be taken for different distances, and the corresponding readings are
noted in the following table.
S.No. Distance Magnet-1 Magnet-2 M1 Tan1
(d)
1.......8 1 1.......8 2 M2 Tan 2
M1
Average
M2 =
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e) Precautions: must be between 300 and 600.
Magnetometer must be gently tapped before noting .
must be noted without parallax error.
Magnetic substances should be removed from the deflection magnetometer.
Method - 2 :
Comparison of Magnetic moments of two bar magnets using a Deflection Magnetometer in tan-A position
by Null Deflection method.
a)Principle: Deflection magnetometer makes use of tangent law.
“ When a magnetic needle is freely suspended in a region where there are two uniform magnetic
fields ‘B’ and ‘BH’ acting at right angles, then needle comes to rest with the axis making an angle
with BH so that B = BH tan ”
b) Description: The deflection magnetometer consists of a small magnetic needle, pivoted on a
point support. A light, thin and long aluminium pointer is attached at right angles to the magnetic
needle. The pointer moves over a scale graduated in degrees as 00 -900 - 00 -900 - 00. A plane
mirror is placed below the scale to observe the deflections without parallax error. The magnetic
needle, pointer and scale with mirror are enclosed in a box with a glass top called compass box. It is
kept at central portion of a wooden bench over which a scale is attached.
c) Arrangement of the magnetometer in tan-A position: The wooden bench is rotated such
that the aluminium pointer becomes parallel to it. i.e., the wooden bench is set in East-West direc-
tion. Without disturbing the bench, the magnetometer is rotated until the aluminium pointer reads 00 -
00. This is called tan-A position.
d) Procedure: One of the magnets of moment M1 is placed on one of the arms at a distance d1.
The second magnet of moment M2 is placed on the second arm and its distance x1 is adjusted such
that the deflection becomes zero. The first magnet is reversed pole to pole in its position. The
second value x2 of the distance of the second magnet is obtained by reversing it pole to pole. Two
more values x3 and x4 for this distance are obtained by exchanging the magnets on the arms and
keeping the first magnet always at the same distance. The average of these four values is noted as d2.
For the deflection to be zero, the magnetic inductions due to the two magnets must be equal and
opposite at the centre of the magnetometer.
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PHYSICS
0 2M 1 0 2M 2 M1 d1 3
. .
4 d13 4 d23 M2 d2 3
A number of observations are taken for different values of ‘ d1’ and are noted in the following
table.
S.No. Distance of Distance of magnet-2 M1 d 1 3
magnet-1 (d1)
x1 x2 x3 x4 Average d2 M2 d 2 3
M1
Average
M2
e) Precautions: Magnetometer must be gently tapped before noting the distance.
All magnetic substances must be removed from the magnetometer.
Method – 3:
Comparison of magnetic moments of the two bar magnets using a deflection magnetometer in tan-B
position by Equal distance method.
a) Principle: Deflection magnetometer makes use of tangent law.
“ When a magnetic needle is freely suspended in a region where there are two uniform magnetic
fields ‘B’ and ‘BH’ acting at right angles, then needle comes to rest with the axis making an angle
with BH so that B = BH tan ”
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IIT - FOUNDATION - SET - V
b) Description: The deflection magnetometer consists of a small magnetic needle, pivoted on a
point support. A light, thin and long aluminium pointer is attached at right angles to the magnetic
needle. The pointer moves over a scale graduated in degrees as 00 -900 - 00 -900 - 00. A plane
mirror is placed below the scale to observe the deflections without parallax error. The magnetic
needle, pointer and scale with mirror are enclosed in a box with a glass top called compass box. It is
kept at the central portion of a wooden bench over which a scale is attached.
c) Arrangement of the deflection magnetometer in tan-B position: The wooden bench is
rotated such that the aluminium pointer becomes perpendicular to it. i.e., the wooden bench is set in
North-South direction, without disturbing the bench, the magnetometer is rotated until the aluminium
pointer reads 00 - 00. This is called tan-B position.
d) Procedure: One of the magnets of moment M1 is placed on one of the arms with its length
perpendicular to the bench. The deflections are noted as 1 and 2 .The magnet is reversed pole to
pole in its position and the deflections are noted as 3 and 4. Four more such deflections are noted
as 5, 6, 7 and 8 keeping the magnet on the other arm and at the same distance. If the average of
these deflections is 1 then,
0 M 1
. BH Tan1 ………(1)
4 d 3
Similarly the average deflection for a second magnet of moment M2 and placed at the same
distance d is obtained as 2 then, 0 . M 2 BH Tan 2 …………….. (2)
4 d 3
M 1 Tan1
From equations and
M 2 Tan 2
A number of observations are taken for different distances and are noted in the following table.
S.No. Distance Magnet-1 Magnet-2
(d) 1..........8 1 1..........8 2
M1
Average =
M2
e) Precautions: must be between 300 and 600.
Magnetometer must be gently tapped before noting q
is to be noted without parallax error.
All magnetic substances must be removed from magnetometer.
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PHYSICS
Method – 4 :
Comparison of Magnetic moments of two bar magnets using a Deflection Magnetometer in tan-B posi-
tion by Null Deflection method:
a) Principle: The Deflection magnetometer makes use of tangent law.
“ When a magnetic needle is freely suspended in a region where there are two uniform magnetic
fields ‘B’ and ‘BH’ acting at right angles, then needle comes to rest with the axis making an angle
with BH so that B = BH tan ”
b) Description: The deflection magnetometer consists of a small magnetic needle, pivoted on a
point support. A light, thin and long aluminium pointer is attached at right angles to the magnetic
needle. The pointer moves over a scale graduated in degrees as 00 -900 -00 -900 - 00. A plane mirror
is placed below the scale to observe the deflections without parallax error. The magnetic needle,
pointer and scale with mirror are enclosed in a box with a glass top called compass box. It is kept at
the central portion of a wooden bench over which a scale is attached.
c) Arrangement of deflection magnetometer in tan-B position: The wooden bench is rotated
such that the aluminium pointer becomes perpendicular to it. i.e., the wooden bench is set in North-
South direction. Without disturbing the bench, the magnetometer is rotated until the aluminium pointer
reads 00 - 00. This is called tan-B position.
d) Procedure: One of the magnets of moment M1 is placed on one of the arms at a distance d1.
The second magnet of moment M2 is placed on the second arm and its distance x1 is adjusted that the
deflection becomes zero.
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The first magnet is reversed pole to pole in its position. The second value x2 of the
distance of the second magnet is obtained by reversing it pole to pole. Two more values x3 & x4 for
this distance are obtained by exchanging the magnets on the arms and keeping the first magnet
always at the same distance. The average of these four values is noted as d2.
For the deflection to be zero, the magnetic inductions due to the two magnets must be
equal and opposite at the centre of the magnetometer.
0 M 1 0 M 2 M1 d13
. .
4 d13 4 d 2 3 M2 d 2 3
A number of observations are taken for different values of d1 and are noted in the following table.
S.No. Distance of Distance of magnet-2 M1 d1 3
magnet-1 (d1)
x1 x2 x3 x4 Average d2 M2 d23
M1
Average
M2
e) Precautions: Magnetometer must be gently tapped before noting the distance.
All magnetic substances must be removed from the magnetometer.
Verification of Inverse square law: (Gauss method)
A) Theory: From the principle of working of deflection magnetometer, we have
B = BH.tan ……… (1)
If a short magnet arranged at a distance ‘d’ end-on to a magnetometer in tan-A position pro-
duces a deflection A then from the Coulomb’s inverse square law, we know that
0 2 M 0 2 M
B= . . BH .tan A ………(2)
4 d 3 4 d 3
If the same magnet, arranged at the same distance but broad side on produces a deflection B in
tan-B position, then from the Coulomb’s inverse square law ; we get
0 M
B= .
4 d 3
0 M
. BH .tan B ……….(3)
4 d 3
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PHYSICS
tan A
Dividing (2) with (3) we get, 2 = tan ……….. (4)
B
tan A
So verifying the result that tan 2 experimentally amounts to the verification of Coloumb’s inverse
B
square law:
B)Procedure: The deflection magnetometer is arranged in tan-A position, as shown in the diagram.
A short bar magnet is placed on one of the arms at certain distance and the deflections are noted as
1, 2. The magnet is reversed pole to pole in the same position and the deflections are noted as 3
and 4. Four more such deflections are obtained as 5, 6, 7 and 8 keeping the magnet on the other
arm at the same distance. The average of these eight deflections gives A. The experiment is re-
peated keeping the magnet at different distances.
The deflection magnetometer is now arranged in tan-B, position as shown in the figure. The
average deflection B is obtained as explained above, keeping the magnet at the same distances as it
is kept in tan-A position.
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The observations are entered into the following tabular form:
S.No. Distance Tan-A Tan-B tan A
‘d’ 1 .....8 A 1 .....8 B tan B
The values in the last column are found to be equal to ‘2’ within the limits of experimental errors. This
verifies the Colulomb’s inverse square law :
C) Precautions: Values of must lie between 300 and 600.
Deflections are to be noted after gently tapping the magnetometer.
Deflections are to be noted without parallax.
Experimental determination of M and BH :
[M = Magnetic moment of the magnet, BH = Horizontal component of Earth’s magnetic induction]
M
Part-1: To find the ratio using deflection magnetometer:
BH
A) Theory: From the principle of working of the deflection magnetometer,
we have B = BH.tan.
0 2M
So, in tan-A position, . BH .tan
4 d 3
M 4 1 3
. .d .tan
BH 0 2
B) Procedure:
The deflection magnetometer is arranged in tan-A position. The given magnet is arranged on one
of the arms at a distance ‘d’ and the deflections are noted as 1, 2. The magnet is reversed pole to
pole in its position and the deflections are noted as 3, 4. Four more such deflections are obtained
keeping the magnet on the other arm at the same distance as 5, 6 , 7 and 8. The average of these
eight deflections gives .
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PHYSICS
The experiment is repeated a number of times keeping the magnet at different distances and
every time measuring the deflections. The observations are entered into the following tabular form:
S.No. distance Deflections d3tan
(d) 1.............8 average
The average of the values in the last column gives d3 tanq.
Now the ratio is calculated using the formula
M
Let this value of be x (1)
BH
Part – 2: To find the product M.BH using the vibration magnetometer:
A) Theory: When a magnet is suspended with its length horizontally by a thread and is disturbed. Then
I
it executes SHM with a period given byT = 2 MB (2)
H
Where I = Moment of inertia of the magnet.
M = Magnetic moment of the magnet.
BH = horizontal component of Earth’s magnetic field.
B) Description: The vibration magnetometer is simply a wooden box with glass doors. The magnet can
be placed in a stirrup, suspended from the torsion head by an unspun silk thread and set to oscillate.
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IIT - FOUNDATION - SET - V
C) Theory: In the equilibrium or mean position the axis of the magnet in along the Earth’s field BH. But
at the instant when the axis of the magnet makes an angle q with the field BH it experiences a couple
C = M BH.sin
C = M BH. - (when = small)
This couple tends to bring the magnet back to the mean position with an angular acceleration. ‘a’ so that
C=I , where I = Moment of inertia of the magnet.
From these equations we have I.. = M.BH..
M .BH MBH
(or) = . . = Constant. (where constant = )
I I
i.e. Angular acceleration a Angular displacement.
So the oscillations are simple harmonic,
2 I
Period T = or T = 2
cons tan t M .BH .
D) Procedure: First, the magnet is suspended and is allowed to come to rest. It is gently disturbed
using another magnet and set to oscillate with small amplitude. The time taken for 20 oscillations is noted using
a stop clock and the average period is obtained as T.
The moment of inertia of the magnet is calculated as
l 2 b2
I=
m
12 12
Where m = mass of the magnet
l = length and b = breadth of the magnet.
I 4 2 .I
But, T =2 or M.BH =
M .BH T2
Let MBH = y
(7)
From equations (1) and (7) we get, M xy = and BH = y/x
E) Precautions:
All other external magnets and magnetic substance are to be removed.
In the case of deflection magnetometer deflection should lie between 300 and 600
In the case of vibration magnetometer deflection should be small enough.
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PHYSICS
ASSIGNMENT
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IIT - FOUNDATION - SET - V
1. Find the magnetic induction at a point mid way between the poles of a horse - shoe magnet, each pole
being 3A-m and distance between them 6 cm.
1) 2/3 ´10-3 N/A-m 2) 4´10-3 N/A-m
3) 10-3 Tesla 4) 3 ´ 10-3 Tesla
2. If the moment of a magnet is 0.4 A-m2 and the force acting on each pole in a uniform magnetic field of
induction 3.2 ´ 10-5 T is 5.12 ´ 10-5 N then the distance between the poles of the magnet is
1) 16 cm 2) 12 cm 3) 25 cm 4) 12.5 cm
3. Magnetic induction at a point on the axial line of a short bar magnet is B towards East. If the magnet is
turned through 900 in clock wise direction, then magnetic induction at the same point is ( Neglect earth’s
magnetic field)
B B
1) towards East 2) towards West
4 2
B B
3) towards North 4) towards south
2 2
4. Two magnets have the same magnetic moments but have different lengths. If the fields at the same
distance from the centre of the magnets in tan A position are measured, which of the following is correct.
1) Both fields are equal 2) Longer magnet produces stronger field
3) Shorter magnet produces a stronger field 4) Cannot decide
5. Two magnets of moments 4 Am2 and 8 Am2 are placed with their equatorial lines coinciding and at a
distance 40 cm apart and with their unlike poles on the same side. The net magnetic induction at a point
mid way between them is
1) 10-4 T 2) 2 ´ 10-4 T 3) 15 ´ 10-5 T 4) 5 ´ 10-5 T
6. A short bar magnet produces magnetic fields of equal induction at two points, one on the axial line and
the other on the equatorial line. Then ratio of their distances
1) 1 : 21/3 2) 1 : 2 3) 21/3 : 1 4) 1 : 8
7. A wire of magnetic moment M is bent at the middle such that the two parts are at an angle 600, with each
other . The new magnetic moment is
2M M 2M
1) 2) 3) 4) 2M
3 2
8. When two magnets of magnetic moments M and M are inclined to each other with a certain angle, the
resultant magnetic moment is found to be M. Then the angle between the magnetic moments is
1) 600 2) 1200 3) 300 4) 450
9. A magnetised wire is bent into an arc of a circle subtending an angle 600 at its centre. Then its
moment is ‘x’ . If the same wire is bent into an arc of a circle subtending an angle 900 at its centre then
its magnetic moment will be
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PHYSICS
2x x 2 2x 3x
1) 2) 3) 4)
3 3 3 2 2
10. A thin straight strip of length 5 cm and magnetic moment 0.5 Am2 was bent such that there is a gap of 1
cm at its ends. Then the magnetic moment of this will be :
1) 0.1 Am2 2) 0.2 Am2
3) 0.05 Am2 4) None
11. A magnet of length 2L and moment ‘M’ is axially cut into two equal halves ‘P’ and ‘Q’. The piece ‘P’
is bent in the form of semi circle and ‘Q’ is attached to it as shown. Its moment is
M M M(2 ) M
1) 2) 3) 4)
2 2 (2 )
N S
N Q S
P
12. Two short bar magnets with magnetic moments 8 Am2 and 27 Am2 were placed 35 cm apart along their
common axial line with their like poles facing each other. The neutral point is
1) mid way between them 2) 21 cm from weaker magnet
3) 14 cm from weaker magnet 4) 27 cm from weaker magnet
13. A bar magnet is placed with its North pole pointing North. Neutral point is at a distance ‘d’ from the
centre of the magnet . The net magnetic induction at the same distance on the axial line of the magnet is
1) 2 BH 2) 3 BH 3) BH 4) 7 BH
14. A bar magnet is placed with its North pole pointing North. Neutral point is at a distance 20cm from the
centre of the magnet . The net magnetic induction at a distance 10cm on the axial line of the magnet is
1) 3 BH 2) 7 BH 3) 15BH 4) 17 BH
15. A bar magnet is placed with its North pole pointing North. Neutral point is at 12 cm . Another magnet
is now placed on the first magnet, then the neutral point is found to be at 8cm. The ratio of their magnetic
moments is
1) 3 :2 2) 27 : 19 3) 9 :4 4) 9: 5
16. Two magnets of moments 4 Am2 and 8 Am2 are placed with their equatorial lines coinciding and at a
distance 40 cm apart and with their unlike poles on the same side. The net magnetic induction at a point
mid way between them is
1) 10-4 T 2) 2 ´ 10-4 T 3) 15 ´ 10-5 T 4) 5 ´ 10-5 T
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IIT - FOUNDATION - SET - V
17. A short magnet is placed with its north pole pointing North. Neutral point is at P. Earth’s horizontal
magnetic induction is BH. If now the magnet is turned through an angle 900. The resultant magnetic
induction at P would be
1) BH 2) 2 BH 3) 5 BH 4) None
18. A very long magnet is held vertically and with its South pole on a table. A single neutral point is located
on the table on the
1) East of the magnet 2) North of the magnet
3) West of the magnet 4) South of the magnet
19. When a bar magnet is placed perpendicular to a uniform magnetic field, it is acted upon by a couple
1.732 ´ 10-5 Nm. The angle through which the magnet should be turned so that the couple acting on it
becomes 1.5 ´ 10 -5 Nm is
1) 600 2) 450 3) 300 4) None
20. Two bar magnets N1S1 and N2S2 having their magnetic moments in the ratio 3 : 1 are rigidly attached
at right angles to each other and the system is pivoted on a free support under the influence of earth’s
field. In equilibrium the angle made by N1S1 with the earth’s field is
1) 600 2) 450 3) 300 4) 900
21. The couple acting on a bar magnet placed in a uniform magnetic field is maximum. The angle through
which magnet is to be rotated so that couple will be half the maximum is
1) 900 2) 450 3) 300 4) 600
22. A magnet is kept fixed with its length parallel to the magnetic meridian. An identical magnet is parallel to
this such that its centre lies on the perpendicular bisector of both. If the second magnet is free to move
it will have
1) translatory motion only 2) rotational motion only
3) both translatory and rotational motion 4) vibrational motion only
23. A magnetic needle, suspended horizontally by an unspun silk fibre, oscillates in the horizontal plane,
because of a restoring force originating mainly from .............
1) The torsion of the silk fibre
2) The force of gravity
3) The horizontal component of earth’s magnetic field
4) All the above factors.
24. A magnet of moment 2 A-m2 lies with its axis in the magnetic meridian, the work done in rotating it
through 370 is (BH = 4 x 10-5 T)
1) 10-5J 2) 2 x 10-5J 3) 1.6 x 10-5J 4) 3.2 x 10-5J
25. When bar magnet is given an angular displacement of 600 in a uniform magnetic field and released its
angular velocity in the mean position will be
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PHYSICS
MB I MB
1) 2MB 2) 3) 4)
I MB I
26. A bar magnet of magnetic moment 2A-m2 is free to rotate about a vertical axis passing through its
center. The magnet is released from rest from east-west position. Then the KE of the magnet as it takes
N-S position is (BH = 25mT)
1) 25mJ 2) 50mJ 3) 100mJ 4) 12.5mJ
27. The period of oscillation of a magnet at a place is 6sec. When it is remagnetised so that the pole strength
becomes 4 times, the period will be
1) 1 sec 2) 2 sec 3) 4 sec 4) 3 sec
28. A bar magnet of length ‘ ’ breadth ‘b’ mass ‘m’ is suspended horizontally in the earth’s magnetic field,
oscillates with period T. If ‘ ‘, m, b are doubled with pole strength remaining the same , the new period
will be
1) 8T 2) 4T 3) T/2 4) 2T
29. The time period of a freely suspended thin magnet is 4sec. If it is cut into two halves perpendicular to its
length and one part is suspended in the same way, the time period will be
1) 0. 25 s 2) 0. 5 s 3) 4 s 4) 2 s
30. A magnet which is suspended by a thread to act as a vibration magnetometer makes 10 oscillations in
one minute at one place 12 oscillations at another place. The ratio of the horizontal components of the
earths magnetic field at these two places is
1) 25 : 36 2) 36 : 25 3) 5 : 6 4) 6 : 5
31. A vibration magnetometer consists of two identical bar magnets placed one over the other such that they
are mutually perpendicular and bisect each other. The time period of oscillation in the earth’s horizontal
magnetic field is 4s. If one of the magnets is taken away, the time period of oscillation of the other in the
same field will be
1) 2 2 s 2) 4 2 s 3) 29/2 s 4) 27/4 s
32. 300 in Tan A position. When the second
Two magnets placed one over the other produce a deflection ofM
1
magnet is reversed pole to pole deflection becomes 600.Then M is
2
1) 3: 2 2) 2:3 3) 1:2 4) 2:1
33. A bar magnet is placed in end -on position at a certain distance from a deflection magnetometer. The
deflection is found to be 600. The magnet is now cut perpendicular to its length in to three equal parts
and one piece is placed at the same distance , the deflection would now be
1) 600 2) 450 3) 300 4) 00
34. In the following , the most sensitive method to compare the magnetic moment of two bar magnets with
a deflection magnetometer is
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IIT - FOUNDATION - SET - V
1) Tan A null method 2) Tan A equidistant
3) Tan B equidistant 4) Tan B null method
35. In working with the deflection magnetometer, the proportional error will be minimum when the deflec-
tion is
1) 900 2) 600 3) 300 4) 450
36. A Robinson’s magnet of pole strength 100Am is placed at a distance of 50cm on one of the arms of a
deflection magnetometer with one of the poles vertically above the magnetic needle. If the deflection
produced is 450, the value of BH is
1) 4 ´ 10-6 T 2) 2 ´ 10-5 T 3) 4 ´ 10-4 T 4) 4 ´ 10-5 T
37. The deflection produced by a magnet in tanA position when it is placed 18cm from the centre is 300.
The deflection in tanB position for another magnet of same length which has 16times the pole strength of
the first and placed at 36cm from the centre is
1) 00 2) 300 3) 600 4) 900
38. A magnet of moment M1 at certain distance end on to a deflection magnetometer produces a deflection
of 600. When a second magnet of moment M2 is placed over it the deflection becomes 600 on the other
side. M1/M2 is
1) 1 : 2 2) 2 : 1 3) 1 : 3 4) 3 : 1
39. A deflection magnetometer is arranged in tan A position. When a bar magnet is placed the deflection is
observed to be ‘q’ and period of oscillation of the needle is ‘x’. When the magnet is removed and
magnetic needle is made to oscillate in presence of earth’s magnetic field then its period is ‘y’. Then
x
y
1) Tanq 2) Cos 3) Sin 4) Cos q
40. The error due to imperfect centering of the pivot of the magnetic needle in the compass box in deflection
magnetometer experiment is eliminated by taking the readings
1) at both ends of Al pointer 2) by reversing the magnet
3) by inter changing the arms of the DMM 4) without parallax
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PHYSICS
KEY
1. 1 16. 4 31. 4
2. 3 17. 3 32. 3
3. 3 18. 2 33. 3
4. 2 19. 3 34. 2
5. 4 20. 3 35. 4
6. 3 21. 4 36. 4
7. 2 22. 1 37. 2
8. 2 23. 3 38. 2
9. 3 24. 3 39. 2
10. 1 25. 2 40.
11. 3 26. 2
12. 3 27. 4
13. 2 28. 4
14. 4 29. 4
15. 2 30. 1
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IIT - FOUNDATION - SET - V
5 ALTERNATING CURRENTS
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PHYSICS
CURRENT ELECTRICITY - III
Alternating Currents
The AC generator or Dynamo : (not for IPE)
The A.C generator consists of a rectangular coil PQRS which is kept rotating with an angular
velocity ‘’ in a uniform magnetic field of induction B in the gap between the poles of a magnet. The axis of
rotation is in the plane of coil and perpendicular to the field B . The ends of the coil are connected to slip rings
R1 and R2, which do rotate with the coils. The AC current is drawn from the carbon brushes B1 and B2 which are
kept pressed against the slip rings.
Theory:
N = number of turns in the coil
A = area bounded by the coil
Let the normal to the plane of the coil make an angle with the field B at an instant ‘t’.
Magnetic flux through the coil = NABcos
= NAB cos t.
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IIT - FOUNDATION - SET - V
The emf induced at this instant is given by
d
E=
dt
E = NAB.sin t.
E = E0 . sin t, where E0 = NAB.
This emf ‘E’ is changing harmonically with time with amplitude or peak value of E0 = NABw.
This harmonic variation of E with time is shown in the graph.
2
The period of revolution of the coil is T = .
At the instant, t = 0....... emf E = 0
t= .......... E = E0
2
2
t= ....... E = 0
2
t= ......... E = -Eo
2
2
t= ........... E = 0.
Thus during one revolution (or) during one cycle, the emf E between B1 and B2 increases from zero to a
peak value, E0 decreases to zero, reversed in polarity increases to a peak value (-E0) and decreases to zero.
Thus, the current drawn during the first half of each cycle is in one direction and during the second half it is
just in the opposite direction. This is the reason why we call the current ‘alternate current’ or simply ‘AC’. This
generator is called the AC dynamo or AC generator. The AC generator in a circuit diagram is represented as
96
PHYSICS
AC circuit with a resistance :
When a resistance R alone is connected to the AC generator of emf E = E0.sint, the current in
the circuit at any instant is given by
E E0
I= .sin t, where E0 is called the peak emf
R R
(or) I = I0 sin t, where I0 is called the peak current.
Thus, the current too is changing harmonically ‘in phase’ with the emf as shown in the graph.
(A) RMS (or) effective values of the AC current and emf :
The average current during each cycle is zero because the current is as much +ve in one half
cycle as it is –ve in another half cycle. But electrical energy is consumed.
This energy that is consumed appears as heat.
The effective or rms value of A.C current is that D.C current which would produce the same
amount of heating as that caused by A.C current in one cycle.
If Peak value of A.C = I0,
then Ieff or Irms =
= 0.707 x I0
(B) Mean or average current :
The mean or avarage current is calculated considering only one half-cycle of A.C current.
The average A.C current is equal to that D.C current which would carry the same amount of charge as
that carried in one half cycle of A.C current.
2
Iave = xI 0
Iave = 0.637 I0
The effective or rms value of the emf is given by
E0
Eeff = Erms =
2
Ex : The household AC supply has Erms = 220V
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IIT - FOUNDATION - SET - V
Hence the peak emf = E0 .Erms
= 1.414 x 220
= 311V
AC circuit with only Inductance : (not for IPE)
Let a coil of inductance ‘L’ and having no resistance be connected to an ideal AC source . Let the current
through the coil at an instant be given by.
I = I0 sint ........ (1)
The emf induced across the inductance or the P.D. across it is just equal to emf of the source and is given
by
dI
E = L.
dt
d
= L. (I sint)
dt 0
= L.. I0 cost.
LI0 represents the peak value E0 of the P.D across the inductor,
thus E = E0. cost.
(or) E = E0 . sin t ......... (2)
2
From equation (1) and (2), it follows that the emf induced across the inductor (i.e., the P.D) is ahead of
the current and the phase difference is .
2
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PHYSICS
The variation of the current ‘I’ and the P.D across the inductor ‘E’ with time is shown in the graph.
At the instant, when the current is zero the e.m.f. has a peak value E0 and at the instant when the emf is zero, the
current has a peak value I0. The resistance offered by the inductance in the AC circuit is called the inductive
reactance.
emf L I 0
Inductive reactance = current I =Lw ohm.
0
If the frequency of AC source is ‘f’ then = 2f.
Inductive reactance XL = L = L.2f
Thus, the reactance offered by an inductor depends on its self inductance ‘L’ and the frequency ‘f’of the
AC current. This inductive reactance increases with increase of frequency.
Note :
The phase difference between the current and the P.D is represented in the form of a vector diagram by
treating the current ‘I’ and the emf ‘E’ as vectors. The emf is p/2 ahead of the current or the current lags behind
the emf by p/2. (w.r.t the anticlockwise rotation)
AC circuit with inductance L and resistance R in series (or) L.R.Circuit (not for IPE)
The AC current in the circuit at any instant is given by
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IIT - FOUNDATION - SET - V
I = I0.sin t........(1)
The P.D. across the resistance at this instant = RI
= RI0 sin wt.
dI
The P.D. across the inductance at this instant = L .
dt
= L. I0.cos t.
= LI0 sin (t + /2)
Vectorially, L I0 sin t is equivalent to LI0 sint oriented at right angles to ‘RI’.
2
This is because, we know that the P.D. across ‘R’ is in phase with the current and the P.D. across
the inductance is /2 ahead of the current.
Vector sum of the potential drops across the R and L at an instant = emf of the source at the same
instant.
emf of the source =
R 2 L2 2 .I
=
R 2 L2 2 I 0 .sin t at an angle ahead of the current I.
E=
R 2 L2 2 .I0 sin (t+) ....... (2)
The quantity is called the “impedance” and is represented by ‘Z’. The emf is ahead of the
current by an angle and is given by
L
tan = ........ (3)
R
L
(or) The current is lagging behind the emf by = tan-1
R
AC circuit with only a condenser : (nor for IPE)
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PHYSICS
A condenser of capacity C is connected to an ideal AC source.
The instantaneous current I is given by
I = I0 sin t. ......... (1)
Q
The P.D. across the condenser =
C
dQ
But the current I = or dQ = I.dt
dt
Q= dQ I dt I .sin t.dt
0
I0
Q= cos t.
The P.D. across the condenser at any instant
is equal to emf of the source at that instant.
Q I
E= 0 .cos t
C C
1
E= .I 0 .sin t ........ (2)
C 2
That is the emf is lagging behind the current by /2. (or) The current is ahead of emf by /2.
The above graph shows the variation of the current I and the emf E with time.
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IIT - FOUNDATION - SET - V
The phase difference between the current and the emf can be represented in the form of a vector diagram
as shown in the figure.
The resistance offered by the condenser in the AC circuit is called the capacitive reactance.
P.D across the condenser
Capacitive reactance =
Current
1
.I 0
1
= C
I0 C
1 1
capacitive reactance XC = ohm.
C C .2 f
The capacitive reactance depends on the capacity of the condenser and the frequency of AC.
the capacitive reactance decreases with increase of the AC frequency.
The C.R circuit (or) AC circuit with a resistance and a condenser.
Let I = I0. sint. be the instantaneous current in the AC circuit containing the resistance ‘R’ and con-
denser of capacity ‘C’.
The instantaneous potential drop across the resistance
(which is in phase with the current ‘I’) is given by
P.D. across ‘R’ = RI = RI0 sin t.
But the instantaneous P.D. across the condenser
1
= I 0 sin t
C 2
1 1
This is equavalent to a P.D. of I o cost or . I oriented with a phase behind the
C. C 2
current.
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PHYSICS
The net P.D. across both R and C represents the emf of the source at the given instant.
This emf is given by
2 1
E = R C 2 2 . I
The phase of this emf is at an angle behind the current ‘I’ (or) The current is ahead of the
emf by angle . This ‘’ is given by
tan =
1
The impedence of this circuit is given by Z = R2 ohm
C 2
2
L.C circuit (or) A.C circuit with pure inductor and capacitor.
Let a pure inductor of inductance ‘L’ and a condenser of capacity ‘C’ be connected in series to the A.C.
source.
Let I = I0 sin t represent the instantaneous current.
At this instant,
a) the P.D. across the inductor is ‘LI’ and is ahead of this current by /2
1
b) the P.D across the capacitor is ‘.I’ and is lagging behind this current by /2
C
1
The net P.D across both L and C L ~ is I and this represents the emf at the same
C
instant.
1
emf E = L ~ .I
C
1
This emf is ahead of the current by /2 if L > and is lagging behind the current by
C
1
/2 if > L.
C
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IIT - FOUNDATION - SET - V
So, the emf is expressed as
1
E = L ~ .I0.sin t .
C 2
Thus the phase difference between the current and the emf in this case is always p/2.
AC circuit with R, L and C in series.
The current in the circuit is I = I0.sint.
where I0 = peak current in the circuit,
= 2f and f = frequency of the source.
The inductive reactance XL = L is /2 ahead of the resistance ‘R’.
The capacitive reactance XC = is /2 behind the resistance ‘R’.
Case – 1 :
If L > then the net reactance is ahead of the
resistance ‘R’ by an angle given by, tan =
2
1 2
and impedence Z = L R
C
The emf ‘E’ is ahead of current by a phase angle and is given by E = E0 sin (t + )
Eo
2
The peak current I0 = 1 2
L R
C
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PHYSICS
Case – 2 :
1
If L then the emf lags behind the current by a phase angle f, given by
C
1
L
tan = C
R
2
1
However impedence Z = L R 2
C
emf E = E0.sin (t-)
now the emf lags behind the current by .
Eo
2
Peak current I0 = 1
L R 2
C
Case-3 :
1
If L = , then they cancel each other. The emf is therefore in phase with the current.
C
E0
Peak current I0 =
R
Thus, for the current in the LCR circuit to be maximum the condition that should be satisfied is
1 1
L = or 2 =
C LC
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IIT - FOUNDATION - SET - V
1 1
(2f)2 = (or) f=
LC 2 LC
This implies that maximum current is drawn from the AC source if the values of ‘L’ , ‘C’ and the
frequency ‘f’ of the AC source satisfy the above condition. This is the case of resonance. By having an AC source
of variable frequency one can study how the current in the circuit changes with ‘f’ for a given set of L and C.
1
Resonant frequency f =
2 LC
Power Factor : (not for IPE)
Power in A. C circuit and the power factor.
The electrical power is the rate at which the electrical energy is consumed and is given by the product EI.
But in an A.C circuit both e.m.f e and the current I keep on changing with time. So the average power’ is to be
calculated considering the r.m.s values of the e.m.f and the current and also the phase difference q between them.
The average power is given by the product of Ir.m.s and the component of Er.m.s in the direction of Ir.m.s
Erms
Irms
Average power = Er.m.s . Ir.m.s cos.
The term ‘cos’ is called the power factor in A.C circuit .
However in an A.C circuit having only a resistance = 0 and cos =1 and
average power = Er.m.s . Ir.m.s
Importance of power factor :
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PHYSICS
With the increase in the reactance the phase q also increases. This is the case when a motor is switched
on. If the inductance of the coil of the motor is ‘L’ then tanq = . Correspondingly the power factor, cosq
decreases and Ir.m.s increases. The transmission power loss in the cable increases abnormally. To minimise such
a loss, a condenser is connected in series when the phase q becomes tanq = .
Then with the inclusion of the condenser the value q is made nearly zero and cosq value is made nearly
one.
When the condenser of your ‘fan’ is burnt, the fan still works but the blades of the fan rotate slowly. The
efficiency of the fan is almost zero though power is consumed. The original working condition can be restored
easily by just replacing the condenser.
Transformer:
The transformer consists of two coils called the primary coil (P) and the secondary coil (S). Both the
coils are wound on the same iron core. Such a transformer works in AC circuits only and is used to convert AC
of one emf into AC of another emf. The coil into which energy is supplied is called the primary coil and the coil
from which the energy is withdrawn is called the secondary.
Theory :
Primary coil :
Number of turns = NP
emf across the coil = EP
current in the coil = IP
Secondary coil :
Number of turns = NS
emf across the coil = ES
current in the coil = IS
Both the coils are wound on the same iron core. Therefore, the rate of change of magnetic flux through
one turn must be same for both the coils. So, if the emf induced across one turn is ‘e’ then
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IIT - FOUNDATION - SET - V
d
e=
dt
This emf ‘e’ per one turn is same for the primary and secondary coils.
emf induced across primary = EP = NP . e
emf induced across secondary = ES = NS . e
Ep NP
........ (1)
ES NS
Also, power input = power output
EP.IP = ES. IS ......... (2)
A) Step up transformer: In this case NS > NP
from equation (1) it follows ES > EP
and from equation (2) it follows IS < IP
So in the case of step-up transformer, low emf and high current are transformed into high emf
and low current.
B) Step down transformer:
In this case NS < NP
from equation (1) and (2) it follows that, ES < EP and IS > IP
So in the case of step down transformer, high emf and low current, are transformed into low emf
and high current.
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PHYSICS
ASSIGNMENT
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1. A coil of area 10-2m2 has 20 turns. This coil is held in a magnetic field of induction 5 tesla and kept
rotating about an axis perpendicular to the field with a frequency of 50Hz. The peak emf induced
between the ends of the coil is nearly
1) 314V 2) 628V 3) 100V 4) 200V
2. The rms value of the emf induced between the ends of the coil in the above problem is nearly
1) 440V 2) 220V 3) 70V 4) 140V
3. A resistance of 100W and an inductance 20mH are connected in series across an ideal AC source. At
the instant when the emf across the source is 120V a current of 1.1A flowing through the resistance. The
rate at which the current through the inductance is changing at that instant is
1) 500A/s 2) 200A/s 3) 1000A/s 4) 400A/s
4. In an AC circuit the inductive reactance and the resistance are in the ration 4 : 3 then
1) the emf lags behind the current by 530 2) the current lags behind the emf by 530
3) the emf lags behind the current by 370 4) the current lags behind the emf by 370
5. In an AC circuit C = 40mF, R = 10W and inductance ‘L’ are in series and the emf is ahead of the current
in phase by 450. If w = 500 rads-1 then L =
1) 20mH 2) 120mH 3) 60mH 4) 1240mH
6. In a series LCR circuit R = 10W , L = 0.1 henry and the instantaneous emf and current are E = 141.4
sin300t and I = I0 sin(300t - 450). The impedance in the circuit is
1) 10W 2) 10 2 W 3) 20W 4) 5 2 W
7. In the above problem the value of I0 is
1) 10A 2) 5A 3) 20A 4) 100A
8. In the problem number 8 the value of C is
1) 166.67mF 2) 0.33mF 3) 33.3mF 4) 0.033mF
9. When 100V D.C is applied across a solenoid a current of 1A flows in it. When 100V A.C is applied
across the same coil the current drops to 0.5A. The frequency of A.C is 50Hz. The impedance and
inductance of the solenoid are
1)100W, 0.75H 2) 100W, 0.60H 3) 200W, 0.55H 4) 200W, 0.75H
10. A 200V-50Hz source is connected across an LCR series circuit. If the potential difference across
inductor and capacitor are 100V each then the potential difference across the resistor is
1) 400V 2)300V 3) 200V 4) zero
11. In an LCR series circuit the voltage across resistance, capacitance and inductance is 10V each. If the
capacitance is short circuited then the voltage across the inductance is
1) 5V 2)10V 3) 10V 4)20V
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PHYSICS
12. A 110V source of alternating current is connected across an LCR series circuit. If resistance of the
circuit is 11W and impedance is 22W then the power consumed is
1)275W 2)366W 3)550W 4)1100W
13. Choose the wrong statement of the following.
1) The peak voltage across the inductor can be greater than the peak voltage of the source in an
LCR circuit
2) In a circuit containing and a capacitor and an ac source the current is zero at the instant source
voltage is maximum.
3) An AC source is connected to a capacitor. The rms current in the circuit gets increased if a
dielectric slab is inserted in to the capacitor.
4) In a pure inductive circuit emf will be in phase with the current.
14. A resistance is connected to an AC source. If a capacitor is included in the series circuit the average
power absorbed by the resistance
1) increases 2) decreases
3) may increase or decrease 4) remains constant
15. A coil, a capacitor and an AC source of voltage 24V are connected in series. By varying the frequency
of the source a maximum rms current of 6A is observed. If this coil is connected to a battery of emf
12V and internal resistance 4Ù the current through it will be
1) 2A 2)1.5A 3)3A 4)2.5A
16. A 120V-60W lamp is run from a 240V-50Hz mains supply using a capacitor connected in series with
the lamp and supply. The capacitor required to operate the lamp at its normal rating is
1) 3.8mF 2) 6.6mF 3) 7.7mF 4)13.3mF
17. At a certain frequency reactance of a capacitor is equal to that of an inductor. If the frequency is
doubled then the ratio of reactance of the inductor to that of the capacitor is
1) 4 : 1 2) : 1 3) 1 : 2 4) 1 : 2
18. A series circuit has an impedance of 50Wand a power factor of 0.63 at 60Hz. The voltage lags the
current. To rise the power factor of the circuit
1) an inductor should be placed in series
2) a capacitor should be placed in series
3) a resistor should be placed in series
4) an inductor or a resistor should be placed in series
19. A filament bulb and an inductor are connected in series to an AC source through a key. The switch is
closed and after some time an iron rod is inserted in to the inductor. Then the glow of the bulb
1) increases 2) decreases 3) remains constant 4) none
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20. A power transmission line feeds input power at 2300V to a step down transformer with its primary
windings having 4000 turns. The number of turns in the secondary in order to get an output power at
230V is
1) 400 2) 4000 3) 40000 4) 92
21. In an oscillating LC circuit the maximum charge on the capacitor is Q. The charge on the capacitor when
the energy is stored equally between the electric and magnetic field is
Q Q
1) Q 2) Q/2 3) 4)
3 2
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PHYSICS
KEY
1. 1 16. 2
2. 2 17. 2
3. 1 18. 4
4. 2
5. 2
6. 2
7. 1
8. 1
9. 3
10. 3
11. 1
12. 1
13. 4
14. 2
15. 2
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6 ELECTRO MAGNETISM
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PHYSICS
CURRENT ELECTRICITY - II
ELECTROMAGNETISM
Oersted’s Experiment:
A straight conductor is held horizontally over a long compass needle NS and parallel to it. When a current
is passed through the conductor, the compass needle is deflected to the position N1S1. On reversing the direction
of the current the needle is found deflected to the other side. This experiment of Oersted showed that the current
carrying conductor is associated with magnetic field.
He also showed that the magnetic lines of force due to a long straight conductor carrying current, in a plane
perpendicular to the conductor are concentric circles.
I I
N1
S
N
S1 I
The direction of the magnetic field associated with a current carrying conductor is given by the following
rules:
(1)Right hand thumb rule : When a conductor carrying a current is held in the right hand so that the
thumb points in the direction of the current, then the curl of the remaining fingers gives the direction of magnetic
lines of induction.
(2)Maxwell’s cork screw rule : If a right handed cork screw is rotated so that its tip advances in the
direction of current, then the direction in which the thumb rotates gives the direction of the magnetic field.
Biot – Savart’s Law or Laplace’s rule :
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IIT - FOUNDATION - SET - V
Statement: The magnitude of magnetic induction ‘dB’ at a point P due to an elementary portion of a
conductor carrying current is directly proportional to
(1) the strength of the current (i).
(2) the length (dl) of the elementary portion
(3) sine of the angle () between the direction of current and line joining the point and midpoint of the
elementary portion.
(4) inversely proportional to the square of the of the distance (r) of the point from the mid point of the
elementary portion.
dB
P
r
dl
I
I
Where K is the proportionality constant whose value depends on the nature of the medium and the system
of units adopted. In S.I. system the unit of dB is tesla.
0
The constant for air or vacuum is K where m0= 4px10-7 henry/metre (Hm-1).
4
0 idl sin
dB . Tesla or NA-1m-1.
4 r2
Magnetic induction at a point near an infinitely long straight conductor carrying current: Let a
current ‘i’ amp. be passing through a long
straight conductor. The magnetic induction B at a point
i
B
r P
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PHYSICS
P at a distance ‘r’ from the conductor is given by
i B
r P
0 2i
B .
4 r
The magnetic induction at a point at a distance r form one end of an infinitely long straight conductor
carrying current is given by
0 i
B .
4 r
Ampere’s Law : The magnetic induction at a point due to an infinitely long, straight conductor is directly
proportional to the strength of the current and inversely proportional to the distance of the point from the
0 2i
conductor. B .
4 r
Ampere’s circutal law : The work done in carrying a north-pole of 1Am, once round the conductor
carrying a current ‘i’ amp is ‘0i’ joule.
Note: This ‘work’ is independent of the shape of the closed path.
Magnetic induction at the centre of a circular coil carrying current :
Let a current of i ampere be flowing through a circular coil of radius ‘r’ having ‘n’ turns. Let MN be a
small portion of length ‘dl’ for one turn. The magnetic induction at the centre ‘O’ due to this elementary portion
is given by
0 idl sin 900
dB
4
.
r 2 900
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IIT - FOUNDATION - SET - V
0 i.dl
dB
4 r 2
The magnetic induction at O due to the whole coil of ‘n’ turns is given by
0 i.dl
B dB .
4 r 2
0i i
2
= dl 0 2 .2r.n
4r 4r
2ni
B 0 . tesla.
4 r
This magnetic induction is perpendicular to the plane of the coil.
0 2 i
If there is only one turn, then B . tesla
4 r
Let ‘P’ be a point on the axis of this circular coil, at a distance ‘x’ from the centre of the coil.
The magnetic induction at P is given by
This magnetic induction acts along the axis of the coil.
Action of magnetic field on a moving charged particle :
direction of this force is given by Fleming’s left hand rule stated as below.
Y
F q ( V x B)
X
B
V
Z
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PHYSICS
Fleming’s left hand rule : If the first three fingers of the left hand are stretched at right angles to each
other, so that the index finger gives the direction of the magnetic induction, the middle finger gives the direction
of motion of the positively charged particle, then the thumb gives the direction of the force experienced by the
charged particle.
In general if a charged particle of charge +q is moving with a velocity ‘v’ in a direction making an angle
with the direction of the magnetic field of induction ‘B’ the force felt by it is given by
__
F qvB sin F q v x B
B
O X
v
Z
So, this force F is zero when is in the direction of and is maximum equal to q v B when and are at the right
angles.
Let a particle having charge ‘q’ and mass ‘m’ enter at right angles into a uniform magnetic field of induc-
tion ‘B’ with a velocity v. This particle experiences a force F = qvB, which is at right angles to ‘v’ at all instants.
So it gets into a circular path of radius ‘r’ such that centripetal force = centrifugal force.
mv 2 mv
qvB= r
r qB
v
qvB
mv2/
r
Action of Magnetic field on current carrying conductor :
When a current carrying conductor is placed in a uniform magnetic field, it experiences a force.
I l B newton. The direction of this force is given by Fleming’s left hand rule.
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IIT - FOUNDATION - SET - V
Statement: If the first three fingers of the left hand are stretched at right angles to each other so that the
index finger gives the direction of the field, the middle finger gives the direction of the current then the thumb
gives the direction of the force.
In general, if the current carrying conductor is held with its length making at an angle of with the direction
of the magnetic field, then the force experienced by it is given by
___
F I . l. BSin or F I ( xB ) .
Torque on a current loop in a uniform magnetic field :
Let a rectangular coil PQRS carrying a current of I amp. be placed with its plane parallel to the direction
of the magnetic field of induction B wb/m2. Let PQ=RS = b and PS = QR = l.
The vertical sides PS and QR experience forces F = IlB newton in the directions as shown in the
figure. These forces are equal, anti parallel and non collinear.
Hence moment of the couple or torque experienced by the coil is given by ,
IlB.b I. lb .B IAB
where A = lb =area of each turn of the coil.
If there are n turns in the coil, then torque is Newton – metre.
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PHYSICS
1) If the coil be placed in the magnetic field such that the normal to its plane makes an angle with the
direction of the magnetic induction, B, then the torque on the coil is
2) The expression for the couple felt by a magnet in uniform magnetic field is C = MB sinq. Comparing
this with the expression for the torque felt by a coil carrying current, in the magnetic field we conclude that a coil
carrying current acts as a magnetic shell of magnetic moment M = niA
S.I. unit of magnetic moment is selected as Amp. m2.
Moving Coil Galvanometer :
Description: A rectangular coil is suspended between the pole faces of a powerful horse-shoe magnet by
a phosphor –bronze suspension fibre and is held in position by a delicate spring. The pole faces are cylindrical to
provide a Radial field. T1 and T2 are the terminals through which current is passed through the coil. A small mirror
strip is fixed to the coil and the deflection of the coil is measured by lamp and scale arrangement. A soft iron
cylinder is arranged within the coil without touching it.
Theory and principle of working :
If n = number of turns of the coil,
A = Area of each turn = l.b
B = Magnetic induction
I = Current through the coil, then
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IIT - FOUNDATION - SET - V
Force on either vertical side = I l B
Torque = I l B x b
The moment of the couple felt by the one turn of the coil is given by
The torque felt by all the ‘n’ turns () = nIAB. This deflects the coil. When the coil is deflected, a restoring
couple is set up in the suspension fibre. If q is the deflection then the restoring couple = c q, where c is called
the restoring couple per unit twist.
In equilibrium position, the deflecting couple = restoring couple.
c.
nIAB c or I = nAB K
c
where K = is a constant called “ figure of merit” of the galvanometer..
nAB
Smaller the value of the “figure of merit”, larger is the sensitivity.
Measurement of the angle of deflection by lamp and scale arrangement.
A collimated beam of light is focussed on to the mirror strip and the reflected spot of light is caught on a
scale. When the mirror deflects through an angle along with the coil then the reflected beam gets deflected
through an angle 2. If the shift on the spot of light on the scale is ‘S’ and the distance of the scale from the mirror
S
is ‘D’ then D.2 = S or
2D
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PHYSICS
Precautions : Since it is very sensitive instrument, only small currents are to be
measured.
Deflection should be small enough so that the spot of light is within the
scale.
Ensure a radial magnetic field between the poles.
Table Galvanometer : The moving coil galvanometer is very sensitive and hence cannot be directly
used in the circuit. Also it is not a portable instrument. The suspended type moving coil galvanometer is modified
into a pivoted type table galvanometer so that it becomes less sensitive and portable. The deflection is measured
by the pointer and scale method.
Thus the table galvanometer is a double pivoted moving coil galvanometer. This pivoted galvanometer
can be moved from place to place.
Tangent Galvanometer : (not for IPE)
(A) Description :
The tangent galvanometer consists of a heavy circular ring which is mounted with its plane verti-
cally on a heavy base provided with levelling screws. This coil can be rotated about a vertical axis through its
centre. The ends of the coil are connected to the terminals T1 and T2. A deflection magnetometer is mounted
on vertical support such that its centre coincides with the centre of the coil.
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( B) Principle :
The tangent galvanometer works on the principle of the tangent law, B = BH tan of the deflection
magnetometer with the difference that the field ‘B’ is due to the current in the coil.
(C) Working :
All magnets and magnetic substances are removed from the place of the experiment. The coil is rotated
slowly until the magnetic needle in the magnetometer is in its plane. That is, the plane of the coil is set in the
magnetic meridian. The magnetometer alone is rotated till the pointer reads 00-00. Now, the current ‘i’ which is
to be measured is passed through the coil. This produces a magnetic field ‘B’ at the centre of the coil in a direction
perpendicular to the plane of the coil. The deflections are noted as 1 and 2. This current is reversed in direction
and the deflections are noted as 3 and 4 . The average deflection is calculated as .
If n = number of turns and r = radius of the coil, then,
o 2 .ni 0 ni
B= . .
4 r 2 r
but B = BH . tan
0 ni
. = BH tan
2 r
2r BH
i=
0 n tan
i = K. tan
The constant ‘K’ is called the reduction factor of the tangent galvanometer. BH is the horizontal compo-
nent of the Earth’s magnetic induction. The current ‘i’ can be measured by knowing K and measuring ‘’.
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PHYSICS
(D) Precautions :
1) Plane of the coil must be in magnetic meridian.
2) Both the ends of the pointer should be read.
3) Experiment is to be repeated by reversing the direction of the current.
Comparison of the working of Moving Coil Galvanometer and Tangent Galvanometer.
Moving Coil Tangent
Galvanometer Galvanometer
Moving coil Moving magnet
galvanometer galvanometer.
Coil carrying the test Small magnetic needle
current rotates in the rotates in the magnetic field
magnetic field of the of the coil carrying the test
magnet. current.
i = k. i = k tan
Since i an evenly I tan. So an evenly
divided scale can be used. divided scale cannot be
used.
The constant k does not The constant k depends
change from place to place on the place of the
of the experiment. experiment as its value
depends on BH.
No initial adjustments Initially the plane of the
are required. coil should be set in the
magnetic meridian.
Working is not effected Working is very much
by the presence of effected by such presence of
magnets or magnetic the magnets and magnetic
substances in the substances.
neighbourhood of the
experiment.
It is not portable. It is portable.
Sensitiveness = 10-8 Sensitiveness = 10-
amp/div. 5amp/div
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Shunt and its theory:
A shunt is a small resistance connected parallel to a sensitive galvanometer so as to protect it from any
damage due to high currents.
G = Galvanometer resistance
S = Shunt resistance
I = Main current
Ig = Current through the galvanometer.
Is = Current through the shunt
I = IS + Ig
Also P.D. across G = P.D. across S
or IgG = (I – Ig)S (since IS= I – Ig)
Ig(S+G) = IS
and
Conversion of a Moving Coil Galvanometer into Ammeter :
(a) A Galvanometer cannot be directly used as an Ammeter because (1) it may not withstand the current
and (2) it has finite resistance and there by the current in the circuit decreases when it is arranged in the circuit.
In other words there is an error in the measurement of current and there is the danger of the Galvanometer getting
damaged.
(b) To overcome these difficulties a shunt resistance is connected across the Galvanometer. The value of
the shunt resistance(s) to be connected depends on the range of the current to be measured by the Ammeter.
If Ig= the current which gives full scale deflection in the
Galvanometer
126
PHYSICS
I = maximum current to be measured by the Ammeter,
G = Galvanometer resistance and S = shunt resistance to be connected.
Then I-Ig = the current through the shunt. Also the potential difference must be same for
both the Galvanometer and the shunt. So, (I – Ig)S = Ig. G
I
If I n , then nI g Ig S Ig G
g
G
S=
n 1
G
So , a shunt resistance should be used to make the range ‘n’ times as large. The equivalent
n 1
resistance of the ammeter is very low. The resistance of an ideal ammeter is zero.
Conversion of a Moving Coil Galvanometer into Voltmeter :
(a) A Galvanometer can not be connected as Voltmeter directly across a resistance in the circuit because
(1) it gets damaged when it draws high current and (2) when it draws some current the measured P.D will be
less than the original PD and hence an error in the measurement.
(b) To overcome these difficulties a high resistance (R) is connected in series with the Galvanometer and
the value of R depends on the range of the P.D. to be measured by the Voltmeter.
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IIT - FOUNDATION - SET - V
If a current Ig gives full scale deflection in the Galvanometer and if it is desired to measure P.D. upto V then
the resistance R that should be used in series is given by the relation.
Ig= where G = Galvanometer resistance
The equivalent resistance of the Voltmeter is very high. The resistance of an ideal Voltmeter is infinity.
Note:
If the range of a Voltmeter is 0 –V and its effective resistance is G then its range can be changed to 0-nV
by connecting an additional resistance
r = G(n-1) in series with it.
Force between two straight parallel conductors carrying current :
Let I1 and I2 be the currents passing through two parallel straight conductors X and Y separated by a
distance ‘r’. Now the magnetic field induction at P on Y due to the current I1 is
0 2 I1
B1 = .
4 r
Force on Y of length is F2 = B1 I2 l
0 2 I1 I 2
F2 = . .l
4 r
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PHYSICS
When ‘X’ exerts a force on ‘Y’, then ‘Y’ also exerts an equal and opposite force on ‘X’.
Force on X = F1 = B2 .F1 . l
0 2I 2
= . .I1.
4 r
0 2I1 .I 2
= . .
4 r
0 2 I1 I 2
Force felt by the either conductor per unit length = . Nm-1
4 r
This force is perpendicular to the length of the conductors. These forces are attractive if I1 and I2 are in
the same direction and ‘ repulsive if I1 and I2 are in opposite direction.
Significance : It gives definition of ampere .
Definition of ampere : One ampere is that current which while passing through each of the two parallel
conductors of infinite length and one meter apart in empty space, causes each conductor to experience a force
of 2x 10-7newton per metre length.
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Electromagnetic Induction :
1) Faraday’s Experiments on Electromagnetic Induction :
Faraday observed that when a magnet is moved near the coil as shown in figure, a deflection is
produced in the galvanometer. And this deflection is found increased when the magnet is brought quickly. This
shows that a current is induced in the coil due to the change of magnetic flux through it. The emf responsible to
it is the induced emf.
A similar phenomenon is again noticed when the magnetic flux through the secondary coil is
changed by changing the current in the primary coil using the key K or the rheostat Rh.
Based on these observations Faraday stated the following laws.
2) Statement of the Faraday’s Laws of Electromagnetic Induction :
i) Whenever there is a change of magnetic flux through a coil, an emf is induced between its ends
and an induced current flows round the coil if it is closed.
ii) The induced current lasts only as long as the magnetic flux through the coil is changing.
iii) The induced emf is directly proportional to the rate of change of magnetic flux through the
coil.
d
the e
dt
3) Statement of Lenz’s Law : The direction of the induced current is such that it opposes the cause to
which it is due.
Explanation : Suppose the N-pole of a magnet is moved along the axis of the coil towards its one face.
Then on looking at this face the induced current I should be anticlockwise so that this face of the coil acts as N-
pole and the coil acts as a magnetic shell. Then only, the coil can repel the approaching N-pole of the magnet as
demanded by the lenz’s law.
For the same reason the induced current will be clock-wise when the N-pole of the magnet is moved away
from the coil.
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PHYSICS
d
4) Fleming’s Right Hand Rule and to show that e =
dt
A straight conductor is held in a magnetic field with its length perpendicular to the field. If the conductor
is moved quickly in a direction perpendicular to its length and perpendicular to the field, then an emf is induced
between its ends and a current tends to flow in the conductor. This is because, it amounts to the change of
magnetic flux while the conductor is cutting through the magnetic field.
Statement of Fleming’s Right Hand Rule: When the thumb, the first finger and the middle finger of
the right hand are stretched at right angles so that the first finger gives the direction of the field, the thumb gives the
direction of motion of the conductor then the middle finger gives the direction of the induced current.
Let a conductor of length ‘l’ be held perpendicular to field ‘B’ and kept moving in a direction perpen-
dicular to both ‘l’ and ‘B’. Let the current induced in the conductor be ‘i’. Referring to the figure, due to the
action of magnetic field on the induced current, the conductor should experience a force (i B) in the downward
direction, as given by the ‘left hand rule’. So work is done in moving the conductor across the field. The work
done in moving it up through a distance ‘s’ is given by
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IIT - FOUNDATION - SET - V
work = iB.s (1)
But (s) represents the area covered by the conductor
s.B represents the magnetic flux through the area cut by the conductor work = i
If the induced emf is ‘e’ and the time taken for this process is ‘t’, then the electrical energy produced = e
i t (2).
From work energy theorem, we have e i t = i e = /t
induced emf = rate of change of magnetic flux
d
e=- (3)
dt
So the constant of proportionality in the Lenz’s law has become one. The negative sign is in accordance
with the Lenz’s law.
Also from equations (1) and (2) it follows that eit = iBs
et = Bs
e = B s/t
So the emf induced in the conductor is given by e = Bv (4)
But s/t represents the velocity v with which the conductor is moved.
5) Self Induction : When a circuit having a coil and a cell is closed, a current begins to flow. In the short
time in which this current increases from zero to a steady value, more and more magnetic lines of force appear
through the coil. This amounts to the change of magnetic flux through the coil. So an emf ‘e’ is induced opposing
this change. Such an opposing emf is called the back emf. The net emf during the growth of the current becomes
(E - e). Similarly, when the current in the coil is decreasing the effective emf becomes (E + e1).
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PHYSICS
Usually when the current is switched off the rate of fall of current or the rate of decrease of the magnetic
field will be very large. So the induced emf e1 becomes very large at the breaking of the circuit. This is the
reason why the sparking occurs in the switch at the break of the circuit.
Definition of self induction: Self induction is that phenomenon in which an emf is induced across the coil
due to the change in its own current.
6) Self Inductance : The magnetic flux ‘‘ through a coil is proportional to the current I
passing through it. I = LI
(1)
The constant L is called the coefficient of self induction or simply the self inductance.
d dI
From equation (1) it follows that L.
dt dt
d
but induced emf, e = -
dt
Definition of the coefficient of self induction : The coefficient of self induction of a circuit is equal to
the emf induced in the circuit while the rate of change of current in it is one ampere per sec. S.I. unit of self
inductance is henry.
Definition of Henry: The self inductance of a circuit is said to be one henry if an emf of one volt is
induced in it while the current through it is changing at the rate of one ampere per second.
7) Mutual Induction: Two coils are said to have mutual inductance when any change in the current in
one coil (primary) induces an emf in the other coil (secondary). The magnetic lines of force due to the current in
the primary are also passing through the secondary coil. Any change in the current in the primary coil causes a
change in magnetic flux through the secondary. So an emf is induced in the secondary coil.
8) Mutual Inductance: The magnetic flux through the secondary coil is directly proportional to the
current I in the primary coil.
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I
= M.I.
(1)
The constant M is called the mutual inductance of the pair of coils.
[Note: It is not the property of any single coil].
d dI
From Equation (1), M.
dt dt
dI
e = - M.
dt
dI
where = rate of change of current in the primary coil and
dt
e = emf induced in the secondary coil
Definition of Mutual Inductance: The mutual inductance of a pair of coils is equal to the emf induced
in one coil while the current in the other coil is changing at the rate of one ampere per sec.
S.I. unit of mutual inductance is henry.
The mutual inductance (M) depends on number of turns in the two coils, the nature of the core present
within the coil, the separation between the coils and how one coil is oriented with respect to another coil.
9. The growth and decay of current through an inductor:
A) Growth of current:
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PHYSICS
A resistance ‘R’ and a coil of self inductance ‘L’ are connected in series with a cell of emf ‘E’ through a
key ‘K’, as shown in the diagram.
When the key is closed, the current in the circuit begins to grow. If ‘i’ is the current at an instant and is the
rate of the growth of this current at the same instant, then at this instant,
P.D. across the resistance = Ri
P.D. across the inductance = L
As seen from the circuit, the cell of emf E has to drive the current against the potential drops across both
the resistance R and the inductance L.
E = Ri + L ........ (1)
Solving this equation mathematically, the value of the current ‘i’ at any instant ‘t’ after closing the circuit is
given by
R
i=
E .t
1 e L
........ (2)
R
The constant e is called the exponential and e = 2.718 » 2.7
The variation of the current ‘i’ with time ‘t’ is shown in the graph. The current is seen growing exponentially
with time.
Analysis:
The P.D. across ‘R’ and ‘L’ are both opposing the emf ‘E’ of the cell. But during the initial stages, ‘i’ is
di di
smaller and is larger. So the P.D. across ‘L’ dominates the P.D. across R. Finally becomes zero or the
dt dt
P.D. across L becomes zero.
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In other words, the inductance is very active initially and becomes passive finally in a D.C. circuit. At this
stage, the cell drives a maximum current ‘i0’ against the resistance R alone.
E
i0 =
R
This current ‘i0’ represents the final ‘steady’ current settled in the circuit. The current at any instant is
therefore given by
R
t
i = i0 1 e
L
........ (3)
L
In a time t = the current grows to a value given by
R
1
i = i0 1
e
1
= i0 1
2.7
i = 0.63 i0
L
This, time t = , is the called the time constant.
R
Definition of time constant:
The time constant of an L.R. circuit is the time taken for the current in the circuit to grow from zero to 0.63
times the maximum steady value.
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PHYSICS
B) Decay of current:
When the circuit is switched off, the current begins to fall rapidly from the steady value of i0 to zero.
As the cell in the circuit is not active, E = 0 .
As the current is decreasing the P.D. across L has a polarity opposite to the P.D across R.
di
0 = Ri + L ............ (1)
dt
di
(or) L . = -Ri
dt
The current ‘i’ at an instant ‘t’ after breaking the circuit can be found from equation (1) as
i = io . ........ (2)
The graph shown represents exponential decay of current with time.
Note :
L
When the circuit is broken, the air gap offers a very large resistance. So the time constant becomes
R
di
very small. Hence, the rate of decay of current becomes very large. A very large emf is therefore induced
dt
across ‘L’ causing the appearance of a spark in the air gap of the switch.
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10. Growth and decay of charge on a condenser :
A) Growth of charge (or) Charging a condenser :
A condenser of capacity ‘C’ is connected to a cell of emf ‘E’ through a resistance ‘R’.
When the key ‘K’ is closed, charge begins to flow from the cell to the condenser through the resistance.
This causes a current i = through the resistance. If Q is the charge on the condenser at an instant ‘t’ and is the rate
of flow of charge then, at this instant,
Q
P.D. across the condenser =
C
dQ
P.D. across the resistance = Ri = R
dt
The emf ‘E’ of the cell is balanced by these potential drops.
Q dQ
E= R ............. (1)
C dt
During the initial stages there is a ‘heavy rush’ of charge to the condenser. That is initially Q has smaller
values and has larger values. In other words, P.D. across the resistance is much larger compared to that across
dQ
the condenser. But, as time passes, the value of ‘Q’ increases and decreases. Finally, a charge ‘Q0’ is
dt
' Q0 ' dQ
settled on the condenser to that the P.D across it, that is is just equal to the emf ‘E’. At this stage, just
C dt
becomes zero (or) the P.D. across the ‘R’ becomes zero. From now onwards ‘R’ becomes passive.
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PHYSICS
In view of this, the equation (1) can be written by
Q0 Q dQ
R. ........ (2)
C C dt
The mathematical solution of this equation gives the value of the charge ‘Q’ on the condenser at any given
instant ‘t’ as
t
Q = Q0 1 e CR ........ (3)
This equation suggests that the charge on the condenser grows exponentially with time, as shown in the
graph.
Also from equation (3), it is seen that in a time t =CR the charge grows to a value,
Q = Q0 1 1 = 0.632 Q0
e
This, time t = C R, is called the time constant.
Definition of time constant :
The time constant of a condenser in a circuit is the time in which the charge on it grows from zero to a value
0.632 times the maximum charge that it can hold in that circuit.
B) Decay of charge on the condenser (or) Discharge of a condenser :
The condenser of capacity ‘C’ is initially holding a charge ‘Q0’. The key ‘K’ is closed to discharge this
condenser through a resistance ‘R’. At an instant ‘t’ after closing the key,
charge on the condenser = Q
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IIT - FOUNDATION - SET - V
dQ
rate of decay of charge = .
dt
In the absence of the cell, we have
Q dQ
0= R. ......... (1)
C dt
However, at the instant t = 0 we have Q = Q0.
The mathematical solution of equation (1) is given by
Q = Q0 . ............ (2)
That is, the charge on the condenser decreases exponentially with time, as shown in the graph.
Qo
1) If V0 = is the initial P.D. across the condenser, then the P.D. at an instant ‘t’ seconds later is given
C
by
t
V = V0 .
e CR
2) In a time t = CR
1
Q = Q0 . = 0.368 x Q0.
e
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PHYSICS
ASSIGNMENT
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1. Flux in weber linked with a closed circuit of resistance 10 ohm varies with time t(in sec) according to
the equation = 6t2 – 5t + 1. Magnitude of the induced current at t=0.25s is
1) 1.2A 2) 0.8A 3) 0.6A 4) 0.2A
2. An aeroplane with wing span 50 m is flying horizontally with a speed of 360 km/hr over a place where
the vertical component of the earth’s magnetic field is 210-4 Wb/m2. The potential difference
between the tips of the wings would be:
1) 0.1V 2) 1.0V 3) 0.2V 4) 0.01V
3. The horizontal component of the earth’s magnetic field at a place is 3 ´ 10-4T and the dip is tan-1(4/3). A
metal rod of length 0.25m placed in the north-south position is moved at a constant speed of 10 cm/s
towards the east. The emf induced in the rod will be:
1) zero 2) 1 V 3) 5 V 4) 10V
4. A rod of length 2m has a translational velocity of 0.5 ms-1 in a direction making an angle 300 with its
length. The plane along which the rod is moving is normal to the magnetic field of induction 2 T. The emf
induced between the ends of the rod is
1) 2V 2) 5V 3) 0.2V 4) 1V
5. A metallic square loop ABCD is moving in its own plane with velocity V in a uniform magnetic field
perpendicular to its plane as shown in the figure. An electric field is induced
1) in AD, but not in BC
2) in BC, but not in AD
3) neither in AD nor in BC
4) in both AD and BC
A B
D C
6. A conducting square loop of side L and resistance R moves in its plane with a uniform velocity V as
shown. A magnetic induction B, constant in time and space pointing perpendicular and into the plane of
the loop exists as shown and towards right side B=0. The current induced in the loop as long as it is
coming out of the field is:
1) clockwise 2) anti clock wise
3) 2BVL/R, anti clock wise 4) zero
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PHYSICS
X X X X
X X X X
V
X X X X
X X X X
7. A metal rod moves at a constant velocity in a direction perpendicular to its length. A constant, uniform
magnetic field exists in space in a direction perpendicular to the rod as well as its velocity. Select the
correct statement from the following:
1) The entire rod is at the same potential
2) There is an electric field in the rod
3) The electric potential is highest at the centre of the rod and decreases towards its ends
4) The electric potential is lowest at the centre of the rod and increases towards its ends
8. A metal bar of length 1m falls from rest under the action of gravity remaining horizontal with its ends in
east-west direction. The induced emf in it at the instant when it has fallen for 10s is (BH = 1.7 x 10–
5
T and g = 10ms–2)
1) 2.5mV 2) 3.2 mV 3) 1.7mV 4) 0.5mV
9. The magnetic field of 2 ´ 10-2 tesla acts at right angles to a coil of area 100 cm2 with 50 turns. The
average emf induced in the coil is 0.1V when it is removed from the field in time t. The value of t is :
1) 0.1s 2) 0.01s 3) 1s 4) 20s
10. A rectangular coil having 600 turns and of area 200cm2 is held with its plane perpendicular to a field of
flux density 8 x10-5T. If the flux density is reduced to 3 x 10-5 Wb/m2 in 0.015s the average e.m.f
induced is
1) 0.01V 2) 0.02V 3) 0.03V 4) 0.04V
11. A copper coil of area of cross section 0.05m2 having 1000 turns is placed in a uniform magnetic field of
4 x 10-5 such that flux through the coil is zero. If it is turned through 900 in 0.01s. If resistance of the coil
is 10W then, the induced charge in 0.01s is
1) 2mC 2) 200mC 3) 0.02mC 4) 20mC
12. A rectangular loop of sides 8cm and 2cm having resistance of 1.6W is placed in a magnetic field of
0.3T directed normal to the loop. The magnetic field is gradually reduced at the rate of 0.02T/s. Power
dissipated is
1)1.6 x 10-10W 2)3.2 x 10-10W 3)6.4 x 10-10W 4) 12.8 x 10-10W
13. A coil of 10 turns has dimensions of 10cm x 10cm. It is rotated at the rate of 7 rev/s about an axis in the
plane of the coil and perpendicular to a uniform magnetic field of flux density 0.5 Wb/m2. The maximum
emf induced is
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IIT - FOUNDATION - SET - V
1) 1.8V 2) 2.2V 3) 2.8V 4) 1.2V
14. A cycle wheel with 64 spokes is rotating with N rotations per second at right angles to horizontal
component of magnetic field. The induced e.m.f. generated between its axle and rim is E. If the number
of spokes is reduced to 32 then the value of induced e.m.f. will be
1) E 2) 2E 3) E/2 4) E/4
15. A wire 40cm long bent into a rectangular loop 15 cm x 5 cm is placed perpendicular to a magnetic field
whose flux density is 0.8 Wb/m2. Within 0.5s the loop is changed into a 10 cm square and flux density
increased to 1.4 Wb/m2. The induced e.m.f is
1) 1.6 x 10-2 V 2) 1.6 x 10-3 V 3) 1.6 x 10-4 V 4) 1.6 x 10-5 V
16. When a nonmagnetic metallic strip is moved away from between the poles of a horse-shoe magnet there
is:
1) a force acting on the strip to oppose the motion
2) a force acting on the strip to help the motion
3) No force acting on the strip
4) a couple acting on the strip so as to rotate it
17. A conducting ring R is placed on the axis of a bar magnet M. The plane of R is perpendicular to this axis.
M can move along this axis then
1) M will repel R when it is moving towards R
2) M will attract R when it is moving towards R
3) M will repel R when moving towards as well as away from R
4) M will attract R when moving towards as well as away from R
R
M
18. A wire AB that lies in the same plane as a circular loop of conducting wire as shown in the figure carries
current from A to B. If current is increasing then what is the direction of the current induced ifany in the
loop?
1) No current will be induced
2) The current will be clockwise
3) The current will be anti-clock wise
4) none
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PHYSICS
A B
19. Two identical coaxial circular loops carry a current I each, circulating in the same direction. If the loops
approach each other, you will observe that:
1) The current in each increases
2) The current in each decreases
3) The current in each remains the same
4) The current in one increases where as that in the other decreases
20. A thin semicircular conducting ring of radius R is falling with its plane vertical in a horizontal magnetic
induction B (figure). At the position MNQ the speed of the ring is V. The potential difference developed
across the ring is :
1) zero
2) BVpR2/V and M is at higher potential
3) pRBV and Q is at higher potential
4) 2RBV and Q is at higher potential
X X X X
B
X X N X X
X X X X
V
X X X X
M Q
21. The network shown in the figure is part of a complete circuit. If at a certain instant the current i is 2A and
is increasing at a rate of 103 A/s., then (VA - VB) is
1) 25V 2) 15V 3) -1V 4) 5V
i
A 6 B
5V 8mH
22. The flux linked with coil changes by 10% , the percentage change in the energy stored in the coil is
1) 10% 2) 20% 3) 21% 4) 42%
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23. Two thick rods AB, CD are placed parallel to each other at a distance l. their ends are joined to a
resistance R. A magnetic field of induction B is applied perpendicular to the plane containing the rods. If
the rods are vertical, the terminal uniform velocity of the rod PQ of mass m is given by
mg.R mg.R mg mgl
1) 2) 3) 4)
B 2l 2 Bl BlR BR
24. The maximum mutual inductance of two coils is equal to 6H. One of the coils has self inductance L1 =
4H. The self inductance of the second coil L2 is.
1) 9H 2) 24H 3) 1.5H 4) 0.67H
25. A conducting rod where is a unit vector along positive Z-axis is moving parallel to x-axis in a uniform
magnetic field directed in positive Y-direction. The end P of the rod will become
1) negative 2) positive
3) neutral 4) sometimes positive and sometimes negative
26. If the current has to flow as shown in the figure, the way that the conductor AB is to be moved in the
magnetic field is; parallel to itself
1) from N to S 2) from S to N
3) into the paper 4) out of the paper
27. Two coils X and Y are placed in a circuit such that when a current changes by 2A in coil X the magnetic
flux linked with Y changes by 0.4 Wb. The mutual inductance of the pair is
1) 0.2 H 2) 5 H 3) 0.2 H 4) 0.8 H
28. A coil has a self inductance 10mH. The maximum e.m.f. induced when the instantaneous current is I =
0.2 sin(100t) Amp., is
1) 0.1 volt 2) 0.2 volt 3) 0.4 volt 4) 0.6 volt
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PHYSICS
29. When current in the primary coil changes from 2amp to 0 amp in 10ms the induced emf in the secondary
is 20 volt. The mutual inductance of the pair of coils is
1) 0.1H 2) 0.2H 3) 0.3H 4) 0.4H
30. A small coil of radius r is placed at the centre of a large coil of radius R , where R >> r. The two coils
are coplanar. The mutual inductance between the coils is proportional to
r r2 r2 r
1) 2) 3) 4)
R R R 2
R2
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KEY
1. 4 16. 1
2. 2 17. 1
3. 4 18. 2
4. 4 19. 2
5. 4 20. 4
6. 1 21. 2
7. 2 22. 2
8. 3 23. 1
9. 1 24. 1
10. 4 25. 1
11. 2 26. 3
12. 3 27. 1
13. 2 28. 2
14. 1 29. 1
15. 1 30. 2
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PHYSICS
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7 ELECTRONICS
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PHYSICS
Introduction
Today computers play an important part in our everyday life. The use of computers today in many fields is
the result of research and development carried out in the field of electronics, the branch of physics that deals with
the study of conduction of electrons and the related phenomenona. From their primitive stage to the present
stage, computers have passed through many generations. In the early stages of development a computer was as
big as a huge building. The essential components that were used in the computers in those days are vacuum tubes,
which are bulky; and that is the reason for the huge size of computers. These computers were referred to as the
first generation computers.
Later, some electronic devices called semiconductor diodes and triodes were invented and used on the
circuits of these electronic devices like computers which increased their speed and efficiency. The replacement of
vacuum tubes with semiconductor diodes and triodes had reduced the size of computers as well. These comput-
ers are referred to as the second generation computers. The development of semiconductor technologyled to the
improvisation of new devices and a large number of semiconductor diodes and triodes were arranged in a single
circuit in a miniature form called integrated circuit (generally referred to as an IC or a chip). The computers that
consist of ICs in their circuits are referred to as the third generation computers. The use of ICs in computers
increased their efficiency. The fourth generation computers contain very large scale integrated circuits, i.e. large
number of chips are integrated in a small circuit; and are referred to as microprocessors. Today microprocessors
are used in the electronic circuits of most of the electronic devices.
In the current chapter the fundamentals regarding the electronic conduction and some of the basic compo-
nents that are the building blocks of an electronic circuit are discussed in addition to the fundamentals of comput-
ers.
Thermionic emission – production of electron stream
It is obvious that electric conduction in solids takes place due to the drift of electrons through them. In
liquids electric conduction takes place through motion of ions; in gases the conduction takes place only at a low
pressure. A gas contained in a tube conducts electricity at a low pressure only when electrodes exist in it, similar
to that in a discharge tube, and when the electrodes are connected to a source of electricity.
Thus electrons are emitted from one electrode and pass through the gas towards another electrode, thereby
conducting electricity through the gas. The electrode from which electrons are emitted is called a ‘cathode’ and
the other electrode is called an ‘anode’. Thus, electric conduction in a gas present in a tube takes place at a low
pressure from the cathode to the anode.
It is found that the movement of the electrons from the cathode to the anode takes place even in vacuum.
Thus electronic conduction is possible in a vacuum tube that contains an anode and a cathode. Such vacuum
tubes were used in the circuitry of the first generation computers, mentioned earlier.
The electrons that are emitted from the surface of the cathode in a vacuum tube are called free electrons.
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IIT - FOUNDATION - SET - V
The emission of free electrons from the surface of a cathode is an important phenomenon. The electrons, that
revolve around the nucleus in the outermost orbit of an atom are valence electrons. These valence electrons are
loosely bound to the nucleus and can be made free by supplying a little energy.
Thus on supplying energy to the cathode of the vacuum tube, free electrons are emitted from its surface.
The energy to be supplied to the loosely bound electrons of the cathode to make them free, can be in the form of
heat, light or high energetic ultra violet or X-ray radiation. If the energy supplied is in the form of heat, the
phenomenon of electron emission is called thermionic emission.
Thus thermionic emission is defined as “the phenomenon in which free electrons are emitted from a metallic
surface, by the absorption of heat energy supplied to them”. The electrons thus emitted are called thermions.
The free electrons of the metallic cathode are free in the sense that they can move about from one atom to
another within the metal. When heat is supplied to the cathode, the free electrons at the surface of the metal gain
enough kinetic energy so that they move away from the cathode and are no more bound to the cathode surface.
The amount of heat energy to be supplied to a metallic surface for it to liberate free electrons (thermions)
is different for different metals and depends on their nature. However, for a given metallic surface, there exists a
minimum heat energy required for a thermion to be liberated and this minimum heat energy is called ‘threshold
energy’ or ‘work function’. The threshold energy or work function is usually measured in electron-volts (eV).
The corresponding temperature of the metal is known as ‘threshold temperature’.
Rate of thermionic emission-factors affecting it
Thermionic emission from a metallic surface is measured by its rate, i.e. the number of thermions emitted
per unit time. This rate of thermionic emission depends on the following factors.
1. The rate of the emission of thermions depends on the nature of the metallic surface. The work
function of each metal is different. If the work function is less, the rate of the emission of thermions is
more. Thus the rate of the emission of thermions is inversely proportional to the work function of
metals.
2. The rate of thermionic emission depends on the temperature of the cathode. The higher is the tem-
perature of the surface of the cathode, the more is the rate of emission of thermions. Thus the rate of
thermionic emission is directly proportional to the temperature of the cathode.
3. The rate of thermionic emission depends on the surface area of the cathode. The more is the surface
area of the cathode, the more is the thermionic emission. Thus, the rate of thermionic emission is
directly proportional to the surface area of the cathode.
Keeping in view the above factors on which the rate of thermionic emission depends, the characteristics of
a good thermionic emitter can be asserted as follows.
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PHYSICS
1. A good thermionic emitter should have a low work function. When an emitter has a low work
function, less amount of heat is energy is required to liberate thermions from its surface. Thus, as
more amount of heat is supplied to it, a larger number of thermions are emitted from it, increasing the
rate of thermionic emission.
2. A good thermionic emitter should have a high melting point. If a thermionic emitter has a high melting
point, it can be heated to a high temperature and thus the rate of thermionic emission can be in-
creased.
Following are some of the thermionic emitters generally used.
(a) Tungsten: It has a high melting point of 3665 K. Its work function is 4.52 electron volt (eV) and
starts emitting thermions at a temperature of about 2500 K.
(b) Thoriated tungsten: A tungsten filament on which a mixture of thorium and carbon is coated is
known as thoriated tungsten. The advantage in this emitter is that heat is produced by passing electric
current through tungsten filament whereas thermions are emitted from thorium. The work function of
thorium is 2.6 eV corresponding to a temperature of about 2000 K. If the temperature is increased
beyond 2000 K, the rate of thermionic emission also increases.
(c) Alkali metal oxides: The work function of alkali metal oxides is much less than that compared to
that of the above-mentioned thermionic emitters. It is less than one electron-volt, and corresponds to
a temperature of about 1000 K. Thus a tungsten filament coated with alkali metal oxides like oxides
of barium, strontium, caesium, etc. is an efficient thermionic emitter.
Kinds of thermionic emitters
Thermionic emitters can be broadly classified into two categories. They are, directly heated thermionic
emitters and indirectly heated thermionic emitters.
In the first category of emitters, current is passed through the emitter, thereby its temperature increases.
When the temperature crosses, the temperature corresponding to the work function of the emitter, it starts
emitting thermions. Tungsten filament is an example of directly heated emitters.
In the second category of emitters, heat is supplied to the emitter by a heating element like tungsten and
when the temperature of the emitter crosses the temperature corresponding to its work function, i.e. its threshold
temperature, it starts emitting thermions. Thoriated tungsten and alkali metal oxides fall under this category of
emitters.
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Applications of thermionic emission
The phenomenon of thermionic emission is used widely in some of the basic components used in electric
circuits of electronic devices like radio, cathode ray tube (CRT) or cathode ray oscilloscope (CRO) etc. One
such basic component is a diode valve.
Diode valve
A diode valve is an electronic device that consists of two electrodes and is used to allow the flow of
electrons in a circuit in a particular direction.
There are two types of diode valves, one that contains a directly heated thermionic emitter and the other
that contains an indirectly heated thermionic emitter.
Directly heated diode valve
The following figure (1), shows the schematic representation of a directly heated diode valve used in an
electric circuit. It consists of a vacuum tube in which a tungsten filament is the cathode and is surrounded by an
aluminium cylinder that acts like an anode. Anode is also referred to as a plate. Figure (2) shows the symbolic
representation of a diode valve in an electric circuit. The filament (F) is connected to a low tension battery (LT)
via an ammeter and a variable resistance. The current through the filament can be changed by varying the resis-
tance and thereby the emission of thermions from the cathode can be controlled. The high-tension battery (HT)
in the circuit provides the potential difference between the cathode and the anode so that the emitted thermions
from the cathode move towards the plate (P) and conducts electricity in the circuit. This diode is not efficient as
the work function of tungsten is high and even the number of thermions emitted are not large.
Glass shell
Filament
Anode cylinder or plate
L.T.
+ F F
+ P V.R
+
H.T
Fig. (1) A directly heated diode valve
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PHYSICS
+
P
+ +
A
L.T. V.R.
H.T.
Fig. (2) Symbolic representation of a directly heated diode valve
Indirectly heated diode valve
In an indirectly heated diode valve, the filament producing heat is surrounded by a cathode that emits
thermions, as shown in figure (3). The terminals of the filament are connected to a low tension battery via an
ammeter and a variable resistance. Adjusting the variable resistance regulates the current through the filament and
thus the regulation of thermionic emission from the filament takes place. The plate and the cathode are connected
to a high tension battery. The symbolic representation of an indirectly heated vacuum diode valve is shown in
figure (4). As the work function and threshold temperature of the cathode is less than that of the filament, thermi-
onic emission takes place at a comparatively lower temperature than that of the filament. As the temperature of
the cathode is increased, more thermions are emitted from its surface.
Glass shell
Filament
Anode cylinder or plate
Cathode cylinder coated
with oxides of barium or
strontium
L.T.
H.T. + P C F F +
Electric flow
A
Fig. (3) An indirectly heated diode valve
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+ +
H.T. L.T. V.R.
Fig. (4) Symbolic representation of an indirectly heated diode valve
Space charge
In a vacuum diode valve, there are two circuits, one for the filament and the other for the conduction of
thermions. Out of these two circuits, if only the filament circuit is switched on, the filament gets heated and
produces thermions. The drift of thermions towards the plate depends on the potential difference between the
plate and the cathode. If the potential difference mentioned above is zero or very less, there is no drift of thermi-
ons emitted from the cathode towards the plate and the emitted thermions are accumulated over the cathode in
the form of electron cloud, called space charge. The space charge around the cathode exists as long as the rate
of drift of thermions towards the anode is less than the rate of thermion production.
Cathode Ray Oscilloscope (CRO)
A cathode ray oscilloscope (CRO) or a cathode ray tube (CRT) is a device that converts electrical signals
into visual signals. The device is used to detect the nature and features of electric signals, different wave forms
and also the changes in the potential difference of alternating currents.
A cathode ray oscilloscope essentially consists of an electron gun, a set of deflecting plates and a fluores-
cent screen. Figure (5) represents a cathode ray oscilloscope.
Electron beam
Vertical hold
plates
Horizontal
hold plates
V.R.
Fluorescent screen
C
+ A A
L.T.
F
S
2000 V
V.R.
H.T. + Control grid
for signal
input ~ ~
A.C. A.C.
Fig. (5) Cathode ray oscilloscope
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PHYSICS
The electron gun consists of a tungsten filament connected to a low tension battery via a switch and a
variable resistance. To the front of the filament, a cathode plate coated with barium oxide or strontium oxide is
mounted and it generates thermions. To the front of the cathode a double anode of cylindrical shape is mounted.
The anode contains a fine hole in it. This is provided so as to accelerate the electrons (thermions) emitted from the
cathode. In front of the anode, a control grid is mounted to which the electric signals are fed.
The fine electron beam emitted from the cathode and accelerated through the anode, after passing through
the grid passes through two pairs of deflecting plates that are connected to a high voltage source of alternating
current. One pair of deflecting plates are placed in vertical position and produces deflections in the horizontal
plane. The other set of deflecting plates are positioned horizontally and cause deflections in the vertical plane.
Thus the operation of these pairs of deflecting plates causes the electron beam to scan the whole screen.
The electron beam that passes between the deflecting plates impinges on the fluorescent screen and causes
scintillations. Thus the electron beam scans the screen with a bluish white glow. The intensity of the electron beam
is controlled by the variable resistance in the low tension circuit. The acceleration of the emitted electrons is
controlled by changing the potential difference between the cathode and the anode in the high tension circuit. The
electron beam emerging through the anode is repelled after it passes through the control grid depending on the
electric signal fed to the grid. This causes variation in the number of electrons striking the screen and hence in the
brightness of the light observed on the screen. If the number of electrons striking the screen are more, a white dot
is observed at the place or else, a black dot is observed. This is the basic principle adopted in television picture
tubes and computer monitors.
Modern electronic devices- An introduction
The vacuum diodes and other related devices that were used once in the circuits of electronic devices were
bulk and thus it was inconvenient to shift them from one place to another easily. But the invention of certain
components in the circuits of electronic devices like semiconductor devices, etc. made them compact and thus
created a revolution in the field of electronics, as mentioned earlier. The components that changed the face of
electronic devices are semiconductor devices and the study of these devices constitutes a separate branch of
physics called ‘solid state physics’. The further content of this chapter deals with the study of semiconductors and
their related devices.
Band theory of solids
Based on the internal structure, solids are classified into two categories, crystalline and amorphous. In an
amorphous solid, the inter atomic distance is not constant and also the atoms of the solid do not have a regular
arrangement. Glass is an example of an amorphous solid. In a crystalline solid, the inter atomic distance is
constant and the atoms have a regular arrangement in space. Pure silicon crystal is an example of a crystalline
solid.
It is obvious that electrons in an atom revolve around the nucleus in elliptical orbits. The energy of electrons
in a given orbit is fixed and thus the orbits are referred to as stationary orbits. The energy of the electrons
increases as the radii of their orbits increases. The energy associated with electrons in a given orbit is referred to
as their energy level. In an isolated atom, the energy levels of electrons are discrete and there exists an energy gap
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between two energy levels that is equivalent to the difference in the energy of the two levels. In a crystal, the
atoms are not isolated. They are fixed to their positions by chemical bonds between them. The interatomic
spacing in a crystal is of the order of one angstrom. In the process of forming a crystal, atoms come closer. In this
process, the outermost orbits of the atoms come much closer. They superpose on one another and finally spread
into a band of energy levels within a small space and this group of energy levels is called an ‘energy band’. The
electrons present in an energy band have no fixed energy and their energy varies continuously within the range of
the energy band. Thus an electron of an atom in a given energy band is free to move within any of the orbits of
atoms present in the band.
Energy bands in a silicon crystal
A silicon crystal entirely consists of silicon atoms. The atomic number of silicon is 14 and its electronic
configuration is 1s2 2s2 2p6 3s2 3p2. In its outermost orbit, there are 4 valence electrons, whereas it can accom-
modate a total of 8 electrons. When many silicon atoms are brought together to form a silicon crystal, their
outermost energy levels are superposed to form an energy band. As the inter-atomic distance is decreased, the
energy band corresponding to 3s level is superposed on the energy bands corresponding to 3p level. In this
situation, the electrons present in the 3s energy levels are free to move even in the 3p energy level. When the inter
atomic spacing is still reduced and is equal to the normal interatomic spacing in a crystal (shown as ‘r’ in figure
(6)), the combined energy band is split into two with a gap called energy gap (EG) formed between them.
Consider that there are ‘n’ number of atoms in a crystal, then the 3s and 3p energy levels of all the ‘n’
number of atoms are superposed to form a ‘3s’ energy band and a ‘3p’ energy band. Since there are ‘n’ number
of atoms, ‘2n’ number of electrons occupy the ‘2n’ number of available ‘3s’ energy states and similarly ‘2n’
number of electrons occupy ‘6n’ number of available ‘3p’ energy states. When all the ‘3s’ and ‘3p’ energy states
are superposed, there are ‘4n’ number of electrons which occupy the ‘8n’ number of energy states. At the normal
inter-atomic spacing, the total energy bands are split into two each having ‘4n’ energy states in each energy band,
all the ‘4n’ electrons in the crystal occupy the ‘4n’ states at the lower level and the ‘4n’ states in the upper level
are left vacant. The energy band at the lower level which is occupied by ‘4n’ electrons is called valence band
(V.B) as the electrons occupy those states are valence electrons. When these valence electrons gain energy
equivalent to the energy gap (EG) and transit to the higher energy level, they are free to move with in the crystal
and conduct electricity. Thus the energy band at the higher level is called conduction band (C.B). Since no
electrons exist in the energy gap (EG), it is also called forbidden gap.
3p
EG
3s
r d
Fig. (6) Energy bands in a silicon crystal
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PHYSICS
At absolute zero temperature, conduction band in any solid does not contain electrons and thus the solid
does not conduct electricity at the temperature. At any higher temperature, some electrons exist in the conduction
band to conduct electricity. Thus the temperature of a solid affects its conductivity.
Conductors, semiconductors and insulators
As discussed earlier, solids are classified into three categories with reference to electric conduction. They
are conductors, semiconductors and insulators.
The following table compares conductors, semiconductors and insulators.
Conductors Semiconductors Insulators
Number of free Of the order Of the order Of the order
electrons 1028 per cm3 1017 per cm3 107 per cm3
Energy gap Nil. Both the Of the order 1 eV. For Of the order 6 eV to
conduction band silicon, it is 10 eV.
and valence band 1.1 eV and for
overlap. germanium, it is
0.7 eV.
Electric Highest Lesser Nil or negligible
conductivity
Temperature As the temperature As the temperature of Similar to that of
coefficient of is increased, the a semiconductor is semiconductors, the
resistance collisions of free increased, more resistance of an
electrons increase number of free insulator is
and thus the electric electrons are released decreased, when its
resistance increases. and so its temperature is
Thus conductors conductivity increased and thus
have positive increases. Thus its its also has a
temperature electric resistance negative temperature
coefficient of decreases with coefficient of
resistance. increase in resistance.
temperature. Hence
semiconductors have
negative temperature
coefficient of
resistance.
Band diagrams
VB Valence C.B C.B C.B
band
CB condition E EG EG
E V.B E
band V.B
EG Energy gap V.B
Usage For electric In electronic devices. As electric shock
conduction. proofers.
Examples All metals and Silicon (Si) and Plastics, rubber
alloys Germanium (Ge) wood, glass, ebonite,
PVC, etc.
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Semiconductors
There are two categories of semiconductors, namely pure and impure. Pure semiconductors are those in
which the entire crystal consists of only one category of atoms. For example, a pure silicon crystal consists
entirely of silicon atoms only. Such a pure semiconductor is called an intrinsic semiconductor. The other
category of semiconductors contains some foreign elements in them. Thus they are called impure or extrinsic
semiconductors.
Intrinsic semiconductor
Consider a pure silicon crystal. A silicon atom contains four valence electrons. Thus every silicon atom
makes four covalent bonds with the neighbouring four silicon atoms. At absolute zero temperature, no free
electrons are available. Thus at this temperature, silicon acts as an insulator. The energy gap between its valence
band and conduction band at this temperature is high and is 1.21 eV. As the temperature is increased, some of the
covalent bonds break up releasing free electrons which move from valence band to conduction band. When an
electron is set free from a bond, it forms an empty space. This empty space formed by the release of a free
electron is called a ‘hole’. The hole is assigned a positive charge as it is formed by a released electron, which is
negatively charged, from that place. Thus, the number of free electrons in an intrinsic semiconductor are equal to
the number of holes in it. As extra charge does not exist on the crystal, it is electrically neutral. Figure (7) shows
the movement of an electron-hole pair in a silicon crystal.
Si
Si Si Si Si Si
Si Si
Si Si
Si Si Si Si Si Si
Si Si O Si Si Si Si Si
Si
Fig. 7(a) Fig. 7(b)
Si Si
Si Si Si Si Si
Si
O
Si Si O Si
Si Si Si Si Si
Si Si Si Si Si Si
Si Si
Fig. 7(c) Fig. 7(d)
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PHYSICS
The following figures, (8) and (9) illustrate the conduction in an intrinsic semiconductor due to electrons
and holes. In the figures, electrons are represented by solid dots (·) and holes are represented by circles (o).
When energy equal to the energy gap of the semiconductor is given, an electron breaks the bond and moves from
valence band (V.B) to conduction band (C.B.) creating a hole at location from ‘A’ (shown in figure (8)). Once a
hole is created at ‘A’, another electron that is let free by breaking another bond at ‘B’ moves to ‘A’ and fills the
hole. Thus an electron moves form ‘B’ to ‘A’ whereas the hole shifts from ‘A’ to ‘B’. Thus the electron movement
and the hole movement are opposite in direction. When a potential difference is applied across the terminals of a
semiconductor crystal, the drift of the free electrons and the holes in it are in opposite directions as shown in
figure (9). Thus the drift of holes in the crystal indicates the direction of conventional current and that of electrons
indicates the direction of electronic current. Hence both the electrons and holes contribute to the current and thus
they are referred to as charge carriers. If ‘In’ indicates the hole current and ‘Ie’ indicates the electronic current, the
total current in the crystal is given by I = In + Ie. Since the electrons and the holes are equal in number, In = Ie and
thus I = 2In = 2Ie.
CB CB CB CB
A B A B A
A C D C
VB VB B VB VB
(a) (b) (c) (d)
Fig. (8) Formation of hole and hole movement
Motion of electrons
Motion of holes
B
Fig. (9) Electric conduction in an intrinsic semicounductor
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Extrinsic semiconductor
The electrons and holes in an intrinsic semiconductor are referred to as charge carriers in it. The number of
electrons and holes in an intrinsic semiconductor are not sufficient to have a better electronic conduction through
it and thus these are not used in any of the electronic devices. To increase the number of charge carriers in an
intrinsic semiconductor and thus to increase the conductivity, the atoms of another element are diffused into it,
thus making the crystal impure. The atoms of the other element added to an intrinsic semiconductor are referred
to as impurities, and the process of the addition of the impurities is called doping. When an intrinsic semiconduc-
tor crystal is doped with a foreign element, the crystal thus formed is called an explicit or an extrinsic semiconduc-
tor crystal. The conductivity of an extrinsic semiconductor is more than that of an intrinsic semiconductor and it
depends on the type of impurity doped and percentage by volume of the doped impurity. Depending on the
nature of the doped material, the extrinsic semiconductors are classified into two categories, namely, p-type and
n-type semiconductors.
p-type semiconductor
Consider a silicon crystal in which indium atoms are doped. The valency of indium (In) is three. Thus it can
form three covalent bonds with its neighbouring silicon atoms in the crystal as shown in figure (10). Hence one
more electron is required to complete the tetravalency of the neighbouring silicon atoms. This implies that a hole
exists at an indium atom in the crystal. Now if some energy is given to the crystal, some covalent bonds break and
free electrons are given out to the conduction band from valence band. But the number of holes present are more
than the free electrons in the crystal. Thus the holes are the majority carriers of charge in the crystal. Thus, the
trivalent impurities added to an intrinsic semiconductor crystal provides more holes in the crystal. Since indium
atom in the crystal accepts an electron, such impurities are referred to as ‘acceptor impurities’. The extrinsic
semiconductor formed by doping a trivalent (acceptor) impurity is called a p-type semiconductor, and a large
number of holes (majority carriers) in the valence band of the semiconductor contribute to the electronic conduc-
tion through it. The minority carriers in this type of semiconductor are electrons and the electronic conduction due
to the electrons is much less compared to that of holes.
Si Si Si
O
Si In Si
Si Si Si
Fig. (10) Bonding in a p-type semiconductor crystal
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PHYSICS
n-type semiconductor
A trivalent impurity in an intrinsic semiconductor results in a p-type semiconductor. But instead of a triva-
lent impurity, if a pentavalent impurity is added to a pure silicon semiconductor, it results in a semiconductor in
which more free electrons (negative charge) exist than the number of holes. Such a type of semiconductor crystal
in which the number of electrons are more than that of the holes in it is called an n-type semiconductor crystal.
Consider a silicon semiconductor crystal doped with a pentavalent impurity like arsenic (As) as shown in figure
(11). Four out of five valence electrons of arsenic take part in the bonding with the neighbouring silicon atoms.
The fifth valence electron is loosely held to the atom. If a little energy is supplied, the electrongets detached
from the atom and becomes free. Thus the number of free electrons in this semiconductor crystal are more than
the number of holes and thus it is called an n-type semiconductor crystal. In this crystal, electrons are the majority
carriers and thus they contribute in large to the electric current. Since each pentavalent impurity like an arsenic
atom, donates a free electron to the crystal, the pentavalent impurities are referred to as donor impurities or
simply donors. The number of holes in this type of semiconductor crystal are less than the number of free
electrons in it and thus they are referred to as minority charge carriers in this crystal.
Si Si Si
Si As Si
free electron
Si Si Si
Fig. (11) Bonding in an n-type semiconductor crystal
p-n junction
A p-type or an n-type semiconductor crystal by itself is not preferred for use in circuits of any electronic
device; as they do not individually have much significance. Hence a suitable combination of these two types of
semiconductor crystals are taken, called p-n junction, to use in the circuits.
Formation of p-n junction
The majority current carriers in a p-type semiconductor are holes and those in an n-type semiconductor
are electrons. The number of majority carriers in them depends on the level of the doping of the impurities in them
and the nature of material being doped. To form a p-n junction, a p-type semiconductor crystal is joined to an n-
type semiconductor crystal by some special methods. At the joint, the p-type and n-type crystals are fused and
thus a p-n junction is formed.
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On either side of the junction, the concentration of holes in ‘p’ type semiconductor and electrons in the ‘n’
type semiconductor is higher; due to this, the holes diffuse from p region to n region and electrons from n region
to p region. Thus at the junction a recombination of holes and electrons take place that causes immobile ions at
the junction. Thus a layer that contains immobile ions is formed at the junction.
This layer is depleted of any charge carriers and thus called a ‘depletion layer’. The ‘depletion layer’ acts
like a wall and restricts further diffusion of holes and electrons into either region across the junction. Thus higher
concentration of holes in the p-region and free electrons in the n-region exists, which gives rise to a building of
positive potential in the p-region and a negative potential in the n-region. Thus there exists a potential difference
across the junction and is referred to as ‘junction potential’. Since the junction acts as a barrier to the flow of
charge carriers, the junction potential is also called ‘barrier potential’.
p type n type
Fig. 12(a)
p J n
D
Fig. 12(b)
p-n junction diode
When two terminals are provided to the ‘p’ type and ‘n’ type semiconductor regions of a p-n junction, the
junction can be treated similar to a vacuum diode, discussed earlier. Hence it is called a p-n junction diode. The
symbolic representation of a p-n junction diode is shown in figure (13). In the symbol the horizontal lines repre-
sent the terminals of the diode. The arrow indicates the direction of the motion of holes, i.e. the conventional
current. Thus the ‘p’ terminal acts as the anode and ‘n’ terminal acts as the cathode.
p terminal n terminal
Fig. (13)
Forward bias and reverse bias of a p-n junction
Consider a p-n junction diode connected in a circuit as shown in figure (14). The ‘p’ type of the diode is
connected to the positive terminal of the battery and the ‘n’ type to the negative terminal. If a p-n junction diode
is connected to a battery in this way, it is said to be in a forward bias condition. When a p-n junction is in forward
164
PHYSICS
bias condition, and a potential difference greater than the barrier potential is applied, then the holes are forced to
diffuse into the ‘n’ region and the electrons into the ‘p’ region, thereby conducting electricity. In this forward bias
condition, when the applied potential is increased, the current through the junction diode is also increased. In this
case, the thickness of the depletion layer decreases.
p n
Fig. (14)
If the ‘p’ and ‘n’ regions of a junction diode are connected to the negative and positive terminals of a
battery respectively, as shown in figure (15), the diode is said to be in ‘reverse bias’ condition. In a reverse bias
condition, the ‘n’ region of the diode is connected to the positive terminal of the battery. Thus the electrons in the
‘n’ region of the diode drift towards the positive terminal of the battery. Similarly, the holes in the ‘p’ region drift
towards the negative terminal of the battery. This makes the current in the circuit almost negligible, compared to
the current in the forward bias. The thickness of the depletion layer increases in this situation. Thus it can be
concluded that the p-n junction diode conducts electricity only in forward bias condition and so it acts like a
vacuum diode valve discussed earlier.
n p
Fig. (15)
Thus the compact p-n junction diodes have replaced the bulky vacuum diode valves in the electronic
circuits, revolutionizing the field of electronics.
Uses of a p-n junction diode
1. Since, a p-n junction diode conducts electricity in one direction only, it can be used as an electronic
switch.
2. These diodes are used in rectification circuits, which convert alternating current into direct current.
This application is of great significance as most of the electronic devices like radio, television, etc.
work using direct current.
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3. A certain category of diodes called light emitting diodes (LEDs) emit light when supplied with elec-
tricity. These diodes are used in digital equipment like, calculators, display boards, traffic signals, etc.
Transistor
A transistor is one of the most important component used in the circuits of electronic devices. It was
invented in 1948 by J. Bardeen, W.H. Brattain and William Shockley. A transistor consists of a combination of
two p-n junctions, unlike a single p-n junction in a p-n junction diode. Hence it is also called a junction transistor.
There are two types of transistors, namely, a p-n-p and an n-p-n transistor.
A p-n-p transistor consists of an ‘n’ type semiconductor sandwiched between two ‘p’ type semiconduc-
tors, giving rise to two p-n junctions. The ‘n’ region in between the two ‘p’ regions is thin, is lightly doped and is
called the ‘base’ of the transistor. Out of the two ‘p’ regions on either sides of ‘n’ region, one region is highly
doped and is called ‘emitter’ whereas the other is moderately doped and is called ‘collector’ in figure (16(a)).
Similarly, in an n-p-n transistor, a ‘p’ type semiconductor which is lightly doped is sandwiched between two ‘n’
type semiconductors giving rise to two p-n junctions. The ‘p’ region is called the base. Out of the two remaining
‘n’ regions, one is highly doped and is the ‘emitter’, the other one is moderately doped and is the ‘collector’
(figure (16 (b)).
p n p n p n
E C E C
B B
(a)
(b)
Fig. (16) (a) p n p transistor and (b) n p n transistor
E emitter, B base, C Collector
The symbolic representation of a p-n-p and an n-p-n transistor are as shown in
figure (17 (a, b)).
E C E C
B B
(a) p n p (b) n p n
Fig. (17)
The direction of arrows in the symbols indicate the direction of the flow of holes, i.e. the conventional
electric current. When a transistor is connected in a circuit, of the two p-n junctions, one is in the forward bias
condition and the other is in the reverse bias condition. The junction which is connected in the forward bias
condition offers less resistance to the current flow and the junction which is connected in the reverse bias condi-
tion offers more resistance. Thus a transfer of resistance takes place in this device and hence it is called a ‘transfer
resistor’ or ‘transistor’.
166
PHYSICS
Uses of a transistor
1. An electronic device that amplifies the input current or voltage signals is called an amplifier. A transis-
tor is an important component in an amplifier circuit.
2. Transistors are also used in electronic circuits called oscillators, which are used to produce signals of
different frequencies.
3. Transistors are used in voltage regulators, also called power stabilizers.
4. Compact transistors are used in minute electronic circuits which include large components, also
called integrated circuits (ICs) or ‘chips’. The chips are used in microprocessors of computers,
calculators, etc.
Radio and television – Broadcasting principles
In the broadcasting of radio and television programmes, electromagnetic waves in the frequency range of
radio frequency are utilised. The fact that the induction coils detect the radio waves transmitted over long dis-
tances was discovered by Marconi in 1895, which led to the invention of radio. The fundamental stages in radio
and television communication are the production of messages, the transmission of the produced messages and
the reception of the transmitted messages. Initially, the messages to be transmitted are converted into electrical
signals.
These message signals are superposed over radio frequency waves, that are capable of propagating over
long distances. These radio frequency waves thus are called ‘radio frequency carriers’. The values of the radio
frequency of the transmitting waves is different for different radio and television stations. For radio transmission,
the radio frequency (r.f) is in the range of 300 kHz to 30 MHz and for television it is in the range of 30 MHz to
300 MHz. The process of superposing a message signal on to a radio frequency wave is known as ‘modulation’.
In modulation, there are two types, namely, ‘amplitude modulation’ (AM) and ‘frequency modulation’ (FM). At
the receiving end, i.e. in the radio set or a television set, these modulated waves pass through a circuit called filter
circuit and the output from the circuit is only a message signal. Thus the message waves are separated from the
modulated waves. This process of extracting message signals from a modulated radio frequency carrier waves is
known as ‘demodulation’ or ‘detection’. The following points summarize the principles of radio and television
broadcasting.
1. The sound or/and the images produced during the making of a program are converted into electrical
signals.
2. The electrical signals produced are amplified by using suitable amplifiers. Amplification is the process
of increasing the amplitude of a given signal.
3. The amplified signal is modulated with a suitable radio frequency carrier wave, generated in an
oscillator.
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4. The modulated waves are fed to the transmitting antenna placed on the top of a tower of consider-
able height.
5. Through the transmitting antenna these modulated waves are transmitted in all possible directions in
the space. This procedure is called transmission.
6. The transmitted modulated waves are received at the receiving antenna when struck with it.
7. The received radio frequency modulated waves, set up alternating currents in the receiving antenna
and through the antenna they are sent to the receiver, i.e. a radio or a television.
8. The tuning circuits in the radio or the television select the desired modulated wave out of the huge
number of waves received at the antenna.
9. The selected modulated wave is fed to the detector circuit, which consists of an electronic oscillator
that separates the message signal (audio or video or both) from the radio frequency carrier wave.
(This process is called demodulation or detection.)
10. The separated message signal is passed through an amplifier circuit which enables us to hear the
sound and / or view the image.
Radio broadcasting and reception
The following block diagram shows the basic principles discussed earlier that are involved in radio com-
munication.
microphone oscillator
frequency wave
generated radio
frequency wave
generated audio
amplifier amplifier
modulator
modulated
wave
Transmitting antenna
Receiving antenna
Tuning circuit
Detector/demodulator
signal
audio
amplifier
Speaker
Fig. (18)
Television broadcasting (telecasting) and reception
168
PHYSICS
T.V. Camera Microphone
Signal
Video
Oscillator
Signal
audio
radio frequency
Oscillator
wave
Video amplifier
Frequency wave
Video Signal
amplified
Audio amplifier
radio
audio signal
Amplified
Amplitude Modulator
Frequency modulator
(AM Picture transmitter)
(FM Sound transmitter)
Combined
Signal
Transmitting antenna
Receiving antenna
Television set
Fig. (19)
In radio broadcasting only audio signals are involved whereas in telecasting, video signals are also involved
along with audio signals. This makes the difference between radio and television broadcasting. The production of
a video signal is done through the scanning of the image of an object. In the early stages of telecasting where there
was no technological advancement, scanning was done using a mechanical system called ‘nipkow disc’. Now-a-
days, an improved system called ‘iconoscope’ is used for the purpose. Iconoscope is a video camera that
consists of a cathode ray beam and photo electric cells. The image of an object is focused onto a screen. The
image formed on the screen is divided into a large number of small square portions.
The light incident on each small square portion is converted into an electric signal by photo cells present on
the screen. The strength of the electric signal produced depends on the intensity of light incident on it. These
electrical (video) signals are amplified further, and then used for modulation. The block diagram in figure (19)
shows the telecasting and reception. The modulation done for video signal is amplitude modulation, whereas the
audio signals are modulated using frequency modulation. Video signals are modulated in ‘AM picture transmitter’
and audio signals are modulated in ‘FM sound transmitter. Then the modulated audio and video signals are
combined and fed to the transmitting antenna placed on a tower. Different television stations use different radio
frequency waves for their transmission and these are referred to as ‘channels’. A television antenna receives these
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transmitted waves and passes them to the tuning circuit of the television set, where a desired channel is selected.
Then the selected modulated wave passes through detection circuit where the messages (both audio and video
signals) are extracted and they are further amplified. The audio and video signals are separated. The audio signals
are amplified and sent to the speaker from which the sound is heard. The video signals are fed to the picture tube
where the electric signals are converted into light and the transmitted image is reconstructed. The television set is
also referred to as kinescope.
Computer
It is a device that helps in computing. It is a multitasking device that performs arithmetic and logical opera-
tions as well. It takes the information (data) given for doing the operations, stores them, processes the informa-
tion and gives the result. It can also store the result (called output). The block diagram of a computer is as shown
in figure (20).
CPU
CU
INPUT OUTPUT
MEMORY ALU
Fig. (20) Block diagram of a computer
Input
The information that is fed to the computer is known as input. The devices through which information is fed
are called input devices. A key board, a mouse, a floppy disc, a compact disc are the different kinds of input
devices.
CPU
CPU is the central processing unit. It is the most important part of a computer and thus is generally referred
to as the heart of it. This unit comprises three sub units. They are control unit (C.U), memory unit (M.U) and
Arithmetic Logic Unit (ALU). Control unit controls the entire system and its related processes. The Memory unit
stores the information and Arithmetic Logic unit performs the operations according to the instructions given.
Output
The processed information is the result and is given out to the user. It is called the output. The devices
through which output is obtained are called the output devices. Monitor, printer, floppy, etc. are the output
devices.
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PHYSICS
The CPU of a computer consists of a circuit board which has an utmost importance and is called ‘mother
board’. The mother board consists of the circuits related to all the components of the computer. These circuits
consist of compact components called micro processors. A micro processor consists of integrated circuits (ICs)
and an integrated circuit consists of a large number of transistors, junction diodes and logic circuit.
A computer cannot perceive the information in the form that we understand, as it uses a binary system. The
binary system consists of only two digits, 0 and 1, ‘0’ refers to ‘off’ mode for a diode and ‘1’ refers to ‘on’ mode.
In the ‘off’ mode, the diode does not conduct electricity. In the ‘on’mode, it conducts electricity. Thus the data
given to a computer is perceived by it in the binary code. The information given to a computer is in the form of
either alphabet or numerals or symbols.
These put together are called ‘characters’. Each character is understood by computer in its binary code. A
binary digit, i.e. either ‘0’ or ‘1’ is called a ‘bit’. A group of eight bits is called a ‘byte’, and a group of one or more
bytes is called a ‘word’. Thus the information is stored in the computer in the form of ‘bytes’ and words’. The
digits in the decimal system, the alphabet and also the special characters in the information given to a computer
are converted into a code that is understood by it before processing. One of the most commonly used codes is
Binary Coded Decimal code (referred to as BCD in short). Generally the numbers, letters and symbols are
represented by 8 bit codes in the BCD code. The first four bits are called ‘zone bits’ and the remaining four are
called ‘numeric bits’. Zone bits help to differentiate between a number, an alphabet and a symbol.
For example, the binary equivalents of the digit 1 and the English alphabet A, both are 0001. Thus to
differentiate between the two, zone bits are used. The zone bit of a numeral is ‘0101’ and that of an alphabet is
‘1010’. Thus the BCD code for ‘1’ is ‘01010001’ and for ‘A’ is ‘10100001’.
Once the computer, perceives the data given to it, the next stage is the processing of the data. This is done
based on the instructions given to it. The instructions given to a computer in the form that we understand, how-
ever cannot be perceived by it. So the instructions are also transformed into the BCD code. The set of instruc-
tions given to a computer is called a ‘program’. The program converted into the BCD code is called ‘machine
language’ and it depends on the CPU. Assembler is an example of a machine language. The physical components
of the computer, i.e. the input and output devices and the CPU are referred to as the ‘hardware’ of the computer
and they work on the machine languages.
But writing the program in a machine language is very tedious. Thus the programs are written in a ‘high
level language’ which consists of English words and symbols, thus making it easy to write a program. There are
certain rules to be followed while writing the programs in a high level language, similar to the grammar rules in
English. These rules are called ‘syntax’.
Once a program is written in a high level language, it has to be translated into a machine language for
execution. The computer takes the help of another program called ‘compiler’ for this translation. Thus the com-
piler is a program that converts a high level language into a machine language.
The complier and the programs written in high level language are referred to as the ‘software’ of the
computer. BASIC, FORTRAN, COBOL, PASCAL, C, JAVA are some of the high level languages.
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Applications of computers
Computers have innumerable applications and today they are part of human life. The following are the
various steps to be followed for solving a problem using a computer.
1. Identifying the problem to be solved by a computer and understanding it.
2. Scripting a suitable program for the problem to be solved by the computer.
3. Supplying data to the computer according to the written program.
4. Getting the output from the computer.
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8 SEMICONDUCTOR DEVICES
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PHYSICS
N-type and P-type extrinsic semiconductors :
N-type semiconductors P-type semiconductors
Semiconductors doped with donor Semiconductors doped with acceptor
atoms (Pentavalent atoms) are called N- atoms (trivalent atoms) are called
type semiconductors. P-type semiconductors
Ex: Silicon + Arsenic Ex: Silicon + Aluminium
Majority carriers are electrons Majority carriers are holes
Number of electrons in the conduction Number of holes in valence band is
band is more than the number of holes more than the number of electrons in
in the valency band. the conduction band.
Fermi level is nearer to the conduction Fermi level is nearer to the valence
band. band.
P-N Junction diode (Semiconductor diode):
The P-N Junction or Semiconductor diode can be formed by adding the N-type impurity in one-
half and P-type impurity in the other half during the growth of a crystal of Germanium(Ge) or Silicon(Si). One
important property of such junction is that it allows the electric current in one direction and offers large resistance
for the flow of current in the opposite direction.
Even in the absence of any P.D. electrons keep migrating from N-region to P-region crossing the junc-
tion. Thereby a layer of +ve charge is developed on the N-side of the junction and a layer of -ve charge is
developed on the P-side. At some stage, the P.D. between the layers becomes so large, that the further flow of
charge across the junction stops. This P.D. is called barrier P.D. which is of the order of 0.1V. These two layers
constitute depletion layer and the width of this layer is about 10-3 mm.
Forward bias: When a battery, having an e.m.f. greater than barrier P.D. is connected to a junction diode
with its +ve terminal to P-region, and -ve terminal to N-region, then the holes are urged to move from P to N and
the electrons from N to P. This establishes a sufficiently large current, which increases with the increase of e.m.f.
of the battery. The P-N junction is now said to be forward biased.
P N P N
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Reverse bias : If the -ve terminal of the battery is connected to the P-region and +ve terminal to the N-
region of a P-N Junction then it is said to be reverse biased.
When a P-N Junction is reverse biased current flow is very small, and is due to the minority carriers. This
current is also called leakage current. The reverse bias increases the width of the depletion layer since it urges the
holes in P-region and electrons in the N-region to move away from the junction.
Since the minority carriers are thermally generated the strength of the reverse current increases with the
increase of temperature. The variation of the current through the diode both under forward and reverse bias is
shown in the above graph.
Thus the P-N junction allows the current to flow in one direction only, and hence it acts as a diode valve.
So it is called the junction diode and is represented in circuit as
The arrow indicates the direction in which the conventional current flows.
(A) Half-Wave Rectifier :
Input Output
Graph - 2
Graph - 1
A rectifier is a circuit which converts A.C. to D.C. The A.C. e.m.f. is changed to the desired value using
a transformer and is applied to a load resistance R through the diode. The diode will be in the conducting state
only during one half of each cycle of A.C. current during which ‘P’ is +ve and ‘N’ is -ve. During this half cycle the
current flows from B to A through the load resistance RL Graph-1 shows the applied A.C. current and graph-2
shows the D.C. current through the resistance R. This current through the resistance RL is direct current.
Efficiency of a half wave rectifier: This is the ratio of the output D.C. power to the input A.C. power.
If ri = diode resistance, RL= load resistance then for a half wave rectifier it can be shown that
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PHYSICS
0.406 Rl
the Rectifier efficiency
ri Rl
This efficiency is maximum when ri = 0
Maximum efficiency : = 0.406 or 40.6%.
But the one disadvantage with this is that the current is pulsating and not continuous.
Keeping this defect in view, a new circuit called full wave-rectifier circuit is built as follows.
(B) Full wave rectifier :
P = primary of the transformer.
S = secondary of the transformer.
C.T= the central tap.
RL = Resistancemsp
The A.C mains e.m.f. is stepped down to the required potential using a step-down transformer.
The ends A and B of the secondary are connected to the P-regions of two diodes D1 and D2. The N-regions are
connected together and connected to the central tap (C.T.) through the load resistance RL. The end A or B of the
secondary becomes alternately +ve and -ve. During each half cycle if one diode is conducting, the other is in the
non-conducting state. So there is always current from C to D through the resistance R. Thus full A.C. wave is
rectified as D.C. in this case. Graph-1 shows the applied A.C. current and graph-2 shows the D.C. current
through the resistance RL.
0.812 RL
Rectifier efficiency for a full wave rectifier is given by = r R
i L
Maximum efficiency for a full wave rectifier is given by =0.812 or 81.2%.
Zener Effect:
From the reverse bias characteristics of a diode it is seen that the reverse bias current increases
from a small value to a very large value suddenly for some potential VZ which is remaining steady. This is called
the Zener effect and the diode working under these conditions is called the zener diode. It is represented by the
following symbol in circuits.
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Note: There are two possible causes for such a steep increase in current namely (1) the breakdown of
the covalent bonds and the release of large number of free electrons and (2) the avalanche ionisation due to
collision of the accelerated free electrons with the atoms.
Zener diode as voltage regulator or stabiliser:
A voltage stabiliser allows us to draw a current of variable magnitude under constant potential
difference. The fact, that the zener potential remains steady though the current through the zener diode, is
changing is used in the working of zener diode as voltage regulator.
The current through the diode may be changing either due to the change in the emf of the cell or
due to a change in the load resistance RL but the potential difference across the zener diode remains constant. In
other words, the zener diode allows the load resistance to draw a current of variable magnitude but under
constant potential difference. Thus the zener diode acts as a voltage regulator or stabiliser.
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PHYSICS
The series resistance R protects the diode from any damage due to high current when RL is made
zero.
V = EMF of the cell - remains constant.
VZ = P.D across the diode - remains constant.
V - VZ = P.D across R - remains constant.
V VZ
I= - remains constant.
R
VZ
IL = R = current through load resistance.
L
IZ = I - IL = current through diode.
Note: When the load resistance RL changes then IL changes, IZ changes but VZ remains constant and I
remains constant.
Transistors:
(1) Two kinds of transistors: In a transistor there are three layers of extrinsic semiconductors with two
junctions in between. So there are three terminals. In the case of a PNP transistor a thin layer of N-type
semiconductor is sandwiched between the thick layers of P-type semiconductors. In the case of an NPN
transistor a thin layer of P-type semiconductor is sandwiched between the thick layers of N-type semiconduc-
tors.
The thin layer is called the base (b) and the other layers are called the emitter (e) and the collector(c).
The emitter is heavily doped, the base is lightly doped and the doping of the collector is in between. These
figures show how these transistors are represented in a circuit. The arrow indicates the direction of the conven-
tional current.
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(2) A transistor in a Circuit - Relation between the currents: The following figures show how the
PNP and NPN transistors are arranged in a circuit with common base (CB) configuration.
The emitter to base junction is forward biased and therefore has a low resistance while
the base to collector junction is reverse biased and has a high resistance.
In the PNP transistor, almost all the holes that pass from emitter to base are simply
urged into the collector region because the base is very thin. So the collector current IC becomes very large. At
the same time, very few holes only are retained by the base and contribute to the base current Ib. Similarly in the
case of NPN transistors almost all the electrons that pass from emitter to base are urged into the collector region.
So the collector current Ic becomes very large. The base, being thin can only retain few electrons. So the base
current is very small.
Applying Kirchhoff’s law we have Ie = Ic + Ib. Also, Ie > Ic >> Ib
(3) Characteristics of a transistor with common emitter (CE) :
Rheostats Rh1 and Rh2 act as potential dividers. The potential difference ‘Vbe’ maintained be-
tween the base and the emitter can be changed using Rh1 and can be measured by a voltmeter. The P.D ‘Vce’
maintained between the collector and the emitter can be changed by using the Rheostat Rh2 and can be measured
by a voltmeter. The base current Ib can be measured using a microammeter and the collector IC can be measured
by using a milliammeter.
(a) Input Characteristics: This graph shows the variation of the input base current Ib with the base to
emitter P.D. ‘Vbe’ while ‘Vce’ is kept constant.
The input resistance ri is defined as the ratio of the change in base to emitter P.D. to the
corresponding change in base current when the collector to emitter P.D. is kept constant.
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PHYSICS
Vbe
ri = I
b Vce cons tan t
The value of ri is variable. So it is said to be non-ohmic.
(b) Output characteristics: This graph shows the variation of the output collector current ‘Ic’with the
collector to emitter P.D. ‘Vce’ ,when the base current Ib is kept constant.
The output resistance “r0” is defined as the ratio of change in collector to emitter P.D. to the change in
collector current when the base current is kept constant.
Vce
r0 = I for the linear portion of the graph.
c Ib cons tan t
(c) Transfer characteristics: This graph shows the variation of the output current Ic with the input
current Ib when Vce is kept constant. The transfer ratio b is defined as the ratio of change in output current to the
change in input current when Vce is kept constant.
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I c
= I
b Vce cons tan t
The value of is around 50. This tells us that a small change in Ib produces a large change in Ic. So the
transfer ratio is also called the current amplification.
Note: If the fraction of the emitter current that appears as collector current is a, then
Ic
Ic = Ie. (or) Ie =
but we have, Ie = Ic + Ib
Ic
=I +I
c b
1
Ic 1 = I
b
1 = I
Ic b
1
Ic. = Ib
I c
I b 1
Current amplification =
1
(d) Voltage amplification: It is the ratio of change in the output collector to emitter P.D. to
change in the input base to emitter P.D.
Vce
A=
Vbe
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PHYSICS
Vbe Vce I c
Note: ri = I , r0
I c = and =
I b
b
r0
A =
ri
(e) Power gain = current gain ´ voltage gain = ´ A
(f) Leakage current and temperature change: Even though the base current Ib is made zero
there is collector current due to the presence of minority charge carriers. This current is called the
leakage current. The population of minority carriers depends very much on temperature. So the
common emitter circuit is very much sensitive for the temperature. It is for this reason silicon transis-
tors are preferred to Germanium transistors. In other words the silicon transistor is less sensitive to
temperature changes.
Common Emitter Amplifier (CE):
Resistances R1 and R2 are used to tap the proper potentials from a common battery respectively
for the base and the collector. The input signal is applied through the condenser C1, between the
base and the emitter. The output signal is obtained between the collector and the emitter. Any small
change in the base current due to the input of a weak signal brings out large change (nearly b times)
in the collector current. So the output signal taken from, across the load resistance is an amplified
signal. This is the principle of working of CE amplifier.
Any change in the P.D Vbe due to input signal Vbe causes a change Ib in the current Ib . This
brings out a large change Ic in the current Ic.
I c
But =
Ib (or) Ic = . Ib
This produces large change in the P.D across RL.
This change in P.D. is the out put signal.
Output signal (P.D) Vce = Ic.RL
Vce
Amplification =
Vbe .
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9 RAY OPTICS
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Rectilinear propagation of light:
Light energy travels in straight lines. The straight path along which the light travels is called the ray of light.
Laws of reflection:
1st law: The incident ray (PQ), the reflected ray (QR) and the normal (QN) at the point of
incidence lie in the same plane.
2nd law: Angle of incidence = Angle of reflection
i.e., i r
a) 1) Real object and virtual object: If the incident rays are diverging then the point from which they
diverge, acts as a real object. If the incident rays are converging then the point at which they tend to converge,
acts as a virtual object.
2) Real image and virtual image: If the rays after reflection or refraction are converging, then the
point at which they converge acts as a real image. Instead, if they are diverging then the point from which they
appear to be diverging acts, as a virtual image.
A real image can be formed on a screen and virtual image cannot be formed on a screen.
Refraction:
A ray of light, while passing from one medium into another medium, usually bends at the surface
of separation. This bending of light is called the refraction.
Laws of refraction:
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PHYSICS
1st law: The incident ray (PQ), the refracted ray (QR) and the normal (N1N2) at the point of incidence
lie in the same plane.
2nd law: Snell’s law of refraction: The ratio of the sine of angle of incidence (i) to the sine of angle of
refraction (r) is a constant, for a given pair of media and for a given colour of light. i.e.,
sin i
=
sin r
1: This constant is called the refractive index or index of refraction of the second medium with respect to
the first medium and is represented by . The value of this constant depends on the nature of the two media and
on the colour of the light.
2: If the first medium is vacuum or air then it is called the absolute refractive index of the second medium
and is represented by 2 or simply m.
3: If C0 = velocity of light in vacuum or air
C1 = velocity of light in medium-1
C2 = velocity of light in medium-2
C1 C0 C0
then = C , 1 , 2
2 C1 C2
2 1
4: From note it follows that 1 2 , 1 2 and 1 3 1 2 . 2 3
1 2 1
Uses of Refractive Index ( or n)
1. The refractive index is a characteristic parameter for a given substance, and for a specific wavelength
. Therefore a substance can be identified by determining its .
Ex : If quartz and calcite crystals are given, they can be identified by measuring their ’s.
2. The refractive index is useful in determining the chemical purity of a substance. A change in the value
for the sodium D line proves that the chemical is not pure.
Total internal reflection and critical angle :
While a ray of light is passing from a denser medium into a rarer medium, both reflection and
refraction take place as long as the angle of incidence (i) is less than certain limiting value called the critical angle
(C). This is what we see for the rays like 1, 2 and 3.
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If the angle of incidence (i) is greater than the critical angle (C) then the total internal reflection takes place
and there would be no refraction. The ray-4 is undergoing total internal reflection.
For the ray - 3, the angle of incidence i = C.
Such a ray grazes the surface of separation so that the angle of refraction r = 900.
sin C 1
With respect to this ray, we have D R =
sin 900
sin C (or) R D =
sin C
If the rarer medium is vacuum or air then the angle C is called the absolute critical angle of the denser
1
medium. In such a case, = .
sin C
Ex: Water : = 1.33 or 4/3 .......... C = 490
glass : = 1.5 or 3/2 .......... C = 420
Diamond : = 2.4 .......... C = 200
Note: The critical angle of a given medium depends on the colour or wavelength of the light.
Consequences of total internal reflection :
1) Mirage, looming and the brilliancy of a diamond are all due to the total internal reflection.
2) Total reflection prism: This is just a rectangular isosceles prism made of glass. Since the angle
of incidence at BC is 450 and is greater than the critical angle of glass (420), it undergoes total internal reflection.
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PHYSICS
3) Optical fibre: An optical fibre consists of three coaxial layers. The innermost layer (A) is a hair thin
fibre of transparent material of =1.7. This layer acts as a light guiding core. This core is covered with a thin
layer (B) of transparent material of =1.6. This layer ‘B’ is known as cladding. The outermost layer (C) acts as
a protective layer.
The critical angle for the boundary between A and B is around 70°. For all angles of launching i1 < 20°, the
angle i2 becomes greater than 70°. So all such rays undergo total internal reflection repeatedly at the boundary
between A and B. So light or a laser signal that enters at one end of ‘A’ arrives at the other end without any
appreciable loss of energy.
Uses of optical fibres:
(1) These fibres are hair thin and are flexible. So they are used in modern communication network. An
electrical signal is first converted into a light signal and the signal is passed through the optical fibre. At the
receiving end the light signal is converted back into electrical signal using a photocell. It is claimed that several
thousand simultaneous telephone conversations can be made through such a fibre.
(2) Endoscopy has become an important method of probing the condition of internal organs of a patient,
without any complicated surgery. Any such endoscope consists of a pair of thin optical fibre cables. One fibre
cable is used to illuminate the internal organ. The other fibre cable is used to collect the reflected light and send
it out. Using this outcoming light signal an image is constructed on the monitor.
(3) Laproscopic or keyhole surgery also employs the optical fibres.
(4) An orthopaedic surgeon uses the optical fibres to probe the condition of the joint before he under-
takes a surgery.
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Refraction through lenses :
a) Optic centre: A ray of light passing through the optic centre of the lens goes undeviated.
b) Focus and Focal length:
fig.(i) fig.(ii)
Convex or converging lens Concave or diverging lens
F = Focus (real) F = Focus (virtual)
The distance of the focus (F) from the optic centre (C) is called the focal length (f).
f = FC (in the figures i and ii )
c) Image formation by a convex lens :
1) Real image formation (model diagram)
OJ = Real Object
IG = Real image
2) Virtual image formation (model diagram)
OJ = Real Object
IG = Virtual image
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PHYSICS
d) Image formation by a concave lens :
OJ - Real object
IG - Virtual image
For all positions of the object ‘OJ’, the image ‘IG’ is erect, virtual, diminished and formed on the same
side between the focus and the optic centre.
e) Relation between the object distance (u), image distance (v) and the focal length (f) for any
lens :
1 1 1 uv
or f
u v f uv
Sign convention:
All distances are measured from the optic centre.
Distances of the real objects and real images are taken to be positive. Distances of the virtual
objects and virtual images are taken to be negative.
Focal length (f) is positive for convex lens and negative for concave lens.
f) Focal power (P) of the lens:
Definition: The focal power of the lens in dioptre is given by the reciprocal of its focal length in metre.
It is positive for convex lens and negative for concave lens.
1 100
P=
f (in meter) f in cm
Smaller the focal length, larger is its converging or diverging action.
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g) Equivalent focal length and equivalent focal power of the lens combination:
1) Lenses in contact :
If two lenses of focal length, f1 and f2 are in contact then the combination acts as a single lens of equivalent
focal length ‘f’ given by the relation.
1 1 1 f1 f 2
or f
f f1 f 2 f1 f 2
The equivalent focal power (P) of the combination is given by P = P1 + P2 , where P1 and P2 are the
focal powers of the two lenses used.
2) Lenses separated by a distance :
If two lenses of focal length f1 and f2 are arranged coaxially at a distance ‘d’ apart, then their equivalent
focal length and equivalent power are given by
1 1 1 d f1 f 2
or f
f f1 f 2 f1 f 2 f1 f 2 d
The equivalent focal power (P) of the combination is given by , where P1 and P2 are the focal powers of
the two lenses used.
Lens maker’s formula:
Let R1 and R2 be the radii of curvature of the two surfaces of a thin lens. is the refractive index of the
medium in which the lens is present and is the refractive index of the material of the lens. AB is the incident ray
running parallel to the axis and F is the focus. OF represents the focal length ‘f’ of the lens.
The focal length of the lens is given by the formula.
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PHYSICS
1 2 1 1 1
f 1 R1 R2
This is the lens maker’s formula.
If a lens made of glass of refractive index ‘m’ is present in air medium of refractive index 1, then
1 1 1
( 1)
f R1 R2
Defects of the image formed by a lens:
1. Spherical aberration:
The spherical aberration in a lens is caused by the fact that the focal length of a lens is more for
paraxial rays and less for marginal rays. As a result of this defect, the images formed by a lens are blurred and are
not sharp enough.
Minimisation of spherical aberration:
i) The spherical aberration in the case of a single lens can be minimised by making the lens surfaces such
that their radii of curvature bear certain ratio depending on the value of m. A lens corrected in this way is called
a crossed lens and it is to be arranged such that the more curved surface receives the incident light.
Ex-1: R1:R2 should be 1:6 for =1.5
Ex-2: A planoconvex lens is enough for =1.6
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ii) If the images are formed by a pair of lenses of focal lengths f1 and f2 separated by a distance ‘d’
then for the spherical aberration to be minimum the following condition should be satisfied.
d = f1 - f2
Also, the lens of larger focal length should receive the incident light.
iii) One rough method of reducing the spherical aberration is to cut-off the paraxial rays by
covering the central portion of the lens with a metal or by simply painting it black. Such a cover that is arranged
is technically called the ‘stop’.
2. Chromatic Aberration :
The chromatic aberration is caused by the fact that the focal length of a lens is least for violet and largest
for red. As a result of this defect the images are not only blurred but also coloured.
There are two types of chromatic aberrations.
i) Longitudinal chromatic aberration : The spread in the coloured image along the axis as shown in
the figure with violet end nearer to the lens is known as the longitudinal chromatic aberration.
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PHYSICS
ii) Lateral chromatic aberration : Formation of coloured images of different size with the
violet image of least size and with red image of largest size is called the lateral chromatic aberration.
Minimisation of chromatic aberration :
i) Achromat : A bi-convex lens made of crown glass and a plano-concave lens made of
flint glass are put in contact using an adhesive called Canada balsam. It is known as achromat.
ii) If images are formed by a pair of lenses separated by a distance ‘d’ then for the chromatic aberration to
be minimum, the lenses of focal length f1 and f2 should be separated by such a distance that d = .
f1 f 2
2
Also the lens of larger focal length should receive the incident light.
Eye-pieces :
If the eye-piece of a telescope or microscope consists of a single lens, then it has the following defects.
a) The image formed by it is less bright as it fails to collect all the light from the objective.
b) This image is not uniformly bright, the central portion being brighter than marginal portion.
c) The spherical and chromatic aberrations are not reduced.
In view of these defects several eye-pieces consisting of a pair of lenses are constructed. For example
: A) Ramsden’s eye-piece B) Huygen’s eye-piece C) Kellner’s eye-piece
A) Ramsden’s eye-piece (+ve eye-piece):
The Ramsden’s eye-piece consists of two plano-convex lenses of same focal length (f), which are ar-
ranged coaxially with their curved surfaces facing each other. They are separated by a distance of 2f/3. The
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RAY OPTICS
‘focus’ of the combination is out side and is at a distance from either lens. So this eye-piece is a positive eye-
piece. Cross wires can be arranged in one of the focal planes. So the telescopes and microscopes which are
intended for measurement are only fitted with Ramsden’s eye-piece. Spherical and chromatic aberrations are
reduced to some extent. The equivalent focal length is 3f/4.
B) Huygen’s eye-piece (- ve eye-piece) :
The Huygen’s eye-piece consists of two plano-convex lenses of different focal lengths, with their
curved surfaces facing towards objective. The focal length of eye lens is f, and the focal length of field lens is 3f.
These two lenses are separated by a distance of 2f, so that spherical and chromatic aberrations
are minimised. Cross wires cannot be introduced in this eye-piece as it is a negative eye-piece. i.e. because focal
plane is in between the lenses. The equivalent focal length of Huygen’s eye piece is 3f/2.
C) Comparison between Ramsden’s and Huygen’s Eyepieces :
Ramsden’s eye-piece Huygen’s eye-piece
It is a positive eye-piece. 1 It is a negative eye-piece.
It does not satisfy the condition for 2 It satisfies the condition for
achromatism. achromatism.
Spherical aberration is reduced to 3 Spherical aberration is removed.
some extent.
Instruments which are intended for 4 It is used for clear observation in
measurements are only fitted with microscopes.
the Ramsden’s eye-piece.
Ex: Telescope of the spectrometer, Ex: dissecting microscope.
travelling microscope.
Convex surfaces of the two lenses 5 Both the convex surfaces face the
face each other. objective.
Dispersion of light:
A) Dispersion:
The splitting up of white light into its constituent colours during its refraction from one medium into
another is called dispersion.
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PHYSICS
B) Spectrum:
When a ray of white light is made to incident on one face of the prism and the emergent beam is caught
on the screen, the coloured patch formed on the screen is called the spectrum.
1 : The angle of incidence ‘i’ is same for all photons. The angle of refraction ‘r’ is least for ‘violet’ and
largest for ‘red’. The angle of deviation ‘D’ is largest for ‘violet’ and least for ‘red’.
2 : Refractive index () of the medium is largest w.r.t. violet and least w.r.t. to red.
C0
3 : = , Velocity ‘C0‘ is same for all photons in vacuum (or) air. Velocity ‘C’ in a medium is least
C
for violet and largest for red.
Spectrum:
a) Impure spectrum: An impure spectrum is one in which the different colours overlap and cannot be
distinguished one from another.
b) Pure spectrum: A pure spectrum is one in which the different colours do not overlap and can be
distinguished one from the other.
c) Arrangement for the production of a pure spectrum :
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Conditions to produce a pure spectrum :
1) The slit must be as narrow as possible.
2) The beam of light from the slit must be rendered parallel by placing the slit at the focus of the convex
lens (L1).
3) The parallel beam of light should be incident on a glass prism in its minimum deviation position.
4) The emergent beam must be brought to focus on a screen by using another convex lens (L2).
Result: Pure spectrum is formed in the focal plane of the second lens. This spectrum is real so it can be
formed on a screen and can also be photographed.
Types of spectra:
Spectra formed by substances are broadly classified into two groups.
I) Emission spectra II) Absorption spectra
I) Emission Spectra: The emission spectrum is obtained by directly using the light from any given
source. Emission spectra are further divided into three types
(a) Continuous spectrum (b) Line spectrum (c) Band spectrum
a) Continuous spectrum: Continuous spectrum is one in which all the colours from violet to red are
present. Each colour is gradually shading off into the next one. It is given by light from solids and liquids in the
state of incandescence.
Ex : The filament lamp, burning candle, burning magnesium wire. Hot charcoal, Kerosene lamp,
heated platinum wire etc.,
b) Line spectrum or Atomic spectra or ionic spectrum: An incandescent gas or vapour in its
atomic state gives a line spectrum. The line spectrum consists of number of sharp and bright lines sepa-
rated from each other by dark spaces. It is a discontinuous spectrum. Such a line spectrum of an element
is characteristically it’s own.
Ex : 1)Hydrogen gas in a discharge tube,
2) The sodium vapour lamp emits light of wavelength 5890 and 5896 causing two lines called D-
lines in the spectrum.
3) Mercury vapour lamp etc.,
c) Band spectrum or molecular spectrum: An incandescent gas or vapour in its molecular s t a t e
gives band spectrum. It consists of number of broad bright bands sharply defined at one edge and dif-
fused at the other. Each band has a thick first line called band lead.
Ex: Nitrogen and Cynogen gases, all compounds of carbon and Bunsen flame.
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PHYSICS
Significance:
1. The line spectrum of an element is characteristically its own. This helps us to identify the elements
present in a given sample.
2. The study of the line spectrum of an element helps us to write down the electron configuration of the
element.
3. The band spectrum of a compound helps us to understand the structure of its molecules. This pro-
cess is called the spectro-chemical analysis.
4. The spectroscopic analysis of the light from a distant star helps us to understand its a) chemical
composition, b) temperature and c) motion.
II) Absorption spectra : If light which is supposed to produce continuous spectrum is ini-
tially passed through a transparent material and then a spectrum of this emergent light is formed, then a spec-
trum crossed by a number of dark lines (or) bands will be observed. Such a spectrum is called absorption
spectrum.
Kirchhoff’s and Bunsen’s Law : According to this law, any substance, which emits light of different
colours when heated to incandescence, possess the property of selectively absorbing these colours when cooled.
Note : The absorption line spectrum of an element is also characteristically it’s own.
Solar spectrum and Fraunhofer lines :
The solar spectrum consists of a continuous spectrum crossed by a large number of dark lines.
These dark lines are called Fraunhofer lines.
When white light from the photosphere of the sun at high temperature (106 0C), passes through
chromosphere (60000C) at low temperature then the chromosphere selectively absorbs photons of certain wave-
lengths, resulting in the appearance of dark lines or Fraunhofer lines in the solar spectrum. From the study of
Fraunhofer lines of the solar spectrum it is possible to detect the kind of elements that are present in the chromo-
sphere of the Sun.
Significance of Fraunhofer lines: Helium was first discovered in the Sun’s atmosphere before it could
be identified on the Earth.
Some dark lines present in the solar spectrum are due to the absorption by the atmosphere of the Earth.
These dark lines are called Telluric lines.
Dispersive power of prism :
The dispersive power of a prism indicates the degree of dispersion that could be caused by it.
Definition: Dispersive power of the prism is defined as the ratio of the angle of dispersion caused by the
prism to the mean deviation caused by it.
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Let b = refractive index w.r.t. blue,
r = refractive index w.r.t. red
= refractive index w.r.t. the mean or yellow
If db, dr and d are the respective deviations
caused by the prism then
b r
db = (b - 1)A, dr = (r - 1)A, d = ( - 1)A, where .
2
The angle of dispersion caused by the prism w.r.t. the blue and red = db - dr
= (b - r) A
Mean deviation caused by the prism d = ( - 1)A
d b d r b r A
Dispersive power =
d 1 A
b r
Dispersive power of the prism =
1
Nature of light:
1. Newton’s corpuscular theory:
i) Light consists of small particles called corpuscles.
ii) Every luminous body emits these corpuscles in all directions with high velocities.
iii) Newton could explain the phenomena like reflection, refraction and rectilinear propagation of light on
the basis of this theory.
Drawbacks:
i) According to him the velocity of light in denser medium is greater than the velocity of light in rarer
medium, which was proved to be wrong later.
ii) This theory, could not explain the phenomena like interference, diffraction polarisation and total internal
reflection.
2. Huygen’s wave theory:
i) Light is propagated in the form of waves in a hypothetical medium called ether, which pervades all space.
This ether is supposed to be extremely rarefied and highly elastic.
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PHYSICS
ii) These waves travel in the form of longitudinal waves. But later Fresnel corrected this as transverse in
nature.
iii) This theory could explain the observed phenomena like rectilinear propagation, reflection, refraction,
diffraction, interference and polarisation.
iv) According to this theory velocity of light in denser medium is less than that in rarer medium.
Drawbacks:
i) There is no existence of the medium ‘ether’
ii) The properties of ‘ether’ are contradicting each other.
iii) The nature of the wave disturbance is not known.
iv) It could not explain photoelectric effect, discrete line spectrum and Compton effect.
3. Maxwell’s electromagnetic Theory:
According to this theory, the propagation of light consists of the propagation of the disturbances in the
intensities of electric field and magnetic field.
i) At any given point both the intensities of electric field and magnetic field change simple harmonically
with time but in directions at right angles to each other.
ii) If the propagation of the wave is along x-axis then sine-wave representing the electric field-wave is in
X-Y plane and the sine-wave representing the magnetic field wave is in X-Z plane.
E = Electric Intensity, B = Magnetic intensity, V = Velocity
iii) At any given point the electric intensity ‘E’ and magnetic intensity ‘B’ become maximum
simultaneously and minimum simultaneously.
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iv) E , B V and and are mutually perpendicular
v) The electric and magnetic fields can propagate in vaccum. So the assumption of the existence of ether
medium is not necessary.
vi) This theory can explain reflection, refraction, Interference, diffraction, polarisation and dispersion.
Drawbacks:
This theory cannot explain 1) Discrete line spectrum
2) Photo electric effect
3) Compton effect
4. Planck’s quantum theory:
i) A photon of discrete frequency ‘’’’ and discrete energy ‘h’ is emitted when an electron takes a jump
from a higher energy level to a lower energy level in an atom.
ii) This photon travels with a velocity equals to velocity of light in vacuum or air
If = permeability and = permitivity of vacuum then the velocity of light in vaccum is given by
1
C0 3 x108 ms 1
0 0
iii) This theory explains best the photoelectric effect, Compton effect and discrete line spectrum.
Drawbacks:
This theory cannot explain 1) Interference 2) Diffraction etc.
Light has both particle and wave properties. This is known as dual nature of light.
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ASSIGNMENT
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1. If refractive index of water is 4/3 and the velocity of light in vacuum is 3 ´ 108 ms 1, then the time taken
by light to travel a distance of 450 m in water is
1) 2 ´ 10-6 s 2) 10-6 s 3) 3 ´ 10-6 s 4) 4 ´ 10-6 s
4 9
2. The refractive index is for water and for diamond. The refractive index of diamond w.r.t water is
3 4
1) 2.25 2) 1.33 3) 3 4) 1.69
3. If eo and mo are the permittivity and permeability of the space e and m are corresponding quantities for
the medium, then the refractive index of the medium is
o o oo
1) . 2) 3) 4)
o o o o
4. If refractive indices of glass and water are 3/2 and 4/3 respectively, the thickness of glass plate in which
light takes same time to travel as of 450 cm of water in cm is
1) 300 2) 400 3) 420 4) 360
5. A glass slab of thickness 4cm contains the same number of waves as 5cm of water when the same
monochromatic ray of light traverses both the media. The refractive index of water is 4/3. The refrac-
tive index of glass is
1) 5/3 2) 5/4 3) 16/3 4) 3/2
6. The critical angle is 600 for medium A and 300 for medium B. The refractive index of B w.r.t A is
1) 3 /2 2) 2/ 3 3) 1.7 4) 3
7. Critical angle will be largest when light goes from
1) Glass to air 2) Water to air 3) Diamond to air 4) Glass to water
8. A light ray travels from denser medium to rarer medium with an angle of incidence ‘i’. Angle between
reflected ray and refracted ray is 900. If the angle of refraction is “r”, critical angle is
1) sin-1(tan r) 2) sin-1(tan i) 3) sin-1(tan r) 4) tan-1(sin i)
9. A ray of light enters a rectangular glass slab of refractive index 3 at an angle of incidence 600. It travels
a distance of 5 cm inside the slab and emerges out the slab. The perpendicular distance between the
incident and the emergent rays is
1) 4Ö3 cm 2) 5/2 cm 3) 5 Ö3/2 cm 4) 5 cm
10. A fish looking up through the water sees the outside world contained in a circular horizon. If the
refractive index of water is 4/3 and fish is 5 7 cm below the surface, what is the radius of the circle?
1) 15cm 2) 12´ 3 cm 3) 12 5 cm 4) 12 7 cm
11. The angles of a prism are 300, 600, 900 ( = 3/2). A transparent liquid film is in contact with the longest
face. A ray normally incident on the shortest face gets just internally reflected at the longest face. The
refractive index of the liquid is
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PHYSICS
1) 2 2) 1.299 3) 3 4) 1.6
12. A beam of light consisting of red, green and blue colours is incident on a right angled prism (as shown in
figure) The refractive indices of the material of the prism for red, green and blue wave lengths are 1.39,
1.44 & 1.47 respectively. The prism will
1) Separate part of the red colour from green and blue colours.
2) Separate part of the blue colour from red and green colours.
3) Separate all the three colour from one another.
4) Not separate even partially any colour from the other two colours.
13. A light ray undergoes a deviation of ‘x’ in a small angled glass prism. (g = 1.5). If it is dipped in a liquid
of refractive index 1.2, then deviation is ‘y’. Then x : y is
1) 4 : 1 2) 1 : 4 3) 1: 2 4) 2 : 1
14. A thin prism of glass is placed in air and water successively. (mg = 3/2, mw = 4/3). The ratio of the angle
of deviations produced by the prism for small angle of incidence when placed in air and water is
1) 9 : 8 2) 4 : 3 3) 3 : 4 4) 4 : 1
15. White light is passed through a prism of angle 50. If the refractive indices for the red and blue rays are
1.641 and 1.659 respectively then the angle of dispersion between them is
1) 0.090 2) 100 3) 200 4) 50
16. A prism with refracting angle of 600, gives angle of minimum deviation 530, 510, and 520 for blue, red
and yellow light respectively. The dispersive power of the material of the prism is
1) 0.38 2) 0.048 3) 0.038 4) none
17. A crown glass prism of angle 6.20 is to be combined with a flint glass prism in such a way that the mean
ray passes un-deviated. If the refractive indices of crown and flint glass for yellow light are 1.517 and
1.620, then the angle of flint glass prism is
1) 1.50 2) 6.20 3) 5.170 4) 1.20
18. In the visible region the dispersive powers and mean angular deviations for crown and flint glass
prisms are 1, 2 and d1, d2 respectively. The condition for deviation without dispersion when the two
prisms are combined is
1) 1 d1 +2 d2 = 0 2) 1 d2 +2 d1 = 0
3) [1 d1]2 +[2 d2]2 = 0 4) [1 d2]2 +[2 d1]2 = 0
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19. Among the following double convex lenses of same material with radii of curvatures and the one that
has largest power is
1) R 1 = , R 2 = 10 cm 2) R 1 = 10cm, R 2 = 5cm
3) R 1 = R 2 = 10 cm 4) R 1 = R 2 = 5 cm
20. Radii of curvatures of convex and concave surfaces of a lens of refractive index 1.5 are 3cm and
4cm respectively. Its focal length is
1) 12 cm 2) 6 cm 3) 24 cm 4) 20 cm
21. A drop of water is placed on a glass plate. A double convex lens having radius of curvature of each
surface 20cm is placed on it. The focal length of water lens (mwater = 4/3) is
1) – 0.20m 2) 0.60m 3) – 0.60m 4) 0.20m
22. A combination of two thin lenses with focal lengths f1 and f2 respectively forms an image of a distant
object at distance 60cm when lenses are in contact. On removing the second lens the image is formed at
a distance 20cm. The values f2 is
1) -60cm 2) -30cm 3) -12 cm 4) -15cm
23. A thin symmetric double – convex lens of focal length f is cut into three parts A, B and C
as shown in figure. The focal lengths of the lenses A, B and C respectively will be
1) f, f, f 2) f, 2f, 2f
3) 2f, f, f 4) f/ 2, f/ 4, f/ 4
B C
24. An equi bi-convex lens has a power of 5 dioptre. If it is made of glass of refractive index 1.5, then the
radius of curvature of each surface is
1) 20cm 2) 7.5cm 3) 6.5cm 4) 40cm
25. The diameter of the aperture of a plano convex lens is 12cm and its thickness at the centre is 2mm. If its
refractive index is 1.5, then its focal length is
1) 20.2cm 2) 40.4cm 3) 90.1cm 4) 180.2cm
26. A double convex lens of focal length 30 cm in air is made of glass of refractive index 1.6. When it is
immersed in a liquid the focal length is found to be 126 cm. The refractive index of the liquid is
1) 1.33 2) 1.5 3) 1.4 4) 1.42
27. If f is focal length of a lens in air, the power of the lens (refractive index = m) when immersed in water
(refractive index = o) is
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PHYSICS
1 1 o 1 ( ) 1
1) f 2) f 3) f 4) ( 1) f
1 o
28. Two equi bi convex lenses, each of focal length 20cm, are placed in contact and the space between
them is filled with water. If glass = 1.5 and water = 4/3 then the focal length of the combination is
1) 10cm 2) 20cm 3) 30cm 4) 15cm
29. The refractive indices of a lens for yellow and red colors are 1.54 and 1.51 respectively. Ratio of the
focal powers of the lens, for yellow light to that of red light is
1) 154 : 151 2) 18 : 17 3) 51 : 54 4) 1 : 1
30. A converging crown glass lens has a focal length 27cm for violet rays. Its focal lengths for red rays is
nearly... (V = 1.5 ; R = 1.45)
1) 24 cm 2) 30cm 3) 25 cm 4) 28 cm
31. A plano convex lens of radius of curvature 15cm and refractive index 1.5 has its curved surface silvered.
It now behaves like a concave mirror of focal length
1) 7.5cm 2) 5cm 3) 10cm 4) 15cm
32. An achromatic doublet used as telescope objective is made up of two materials of dispersive powers
0.012 and 0.018. Its effective focal length is 120cm. The focal length of the convex lens is
1) 40cm 2) 50cm 3) 60cm 4) 20cm
33. The dispersive powers of the materials of two lenses forming an achromatic doublet are in the ratio 3 :
4. If the focal length of the combination is +60cm, the focal lengths of component lenses are
1) –20, 25 2) 20, -25 3) –15, 20 4) 15, - 20
34. The equivalent focal length of a Ramsden’s eyepiece is 2.4cm. The distance between the two lenses is
1) 2.13 cm 2) 0.9cm 3) 3.6cm 4) 1.2cm
35. The equivalent focal length of Huygen’s eyepiece with eyelens of focal length f is
1) 3f/4 2) 3f/2 3) 4f 4) 2f
36. An eye piece is constructed so that spherical and chromatic aberrations are minimised to the maximum
extent. If its equivalent focal length is ‘p’ then the focal length of the field lens is
3p 2p
1) 2) 2p 3) 4) 3p
2 3
37. The spectrum emitted by a white hot solid is
1) Discrete line spectrum 2) Discrete band spectrum
3) Continuous spectrum 4) Absorption spectrum
38. When light from white-hot solid passes through sodium vapours, the spectrum of emergent beam will
show
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1) Two bright lines in the yellow region 2) Two dark lines in the yellow region
3) All colours of the rainbow 4) no colours at all
39. The refractive index of the core of an optical fiber is 2 and that of the cladding is m1. The angle of
incidence on the face of the core so that the light ray just undergoes total internal reflection at the
cladding is
1 1
1) sin 2) sin 1 22 12 3) sin 1 2 1 4) sin 1 22 12
2
40. Huygen’s principle is used to find
1) the velocity of light 2) the position of the wave front
3) the wavelength of light 4) the focal length of lens
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PHYSICS
KEY
1 1 16 3 31 2
2 4 17 3 32 1
3 4 18 1 33 4
4 2 19 4 34 1
5 1 20 3 35 2
6 4 21 3 36 2
7 4 22 2 37 3
8 2 23 2 38 2
9 2 24 1 39 2
10 1 25 4 40 2
11 2 26 3
12 1 27 4
13 4 28 4
14 4 29 2
15 1 30 2
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10 PHYSICAL OPTICS
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PHYSICS
Introduction:
Images are formed on a screen by placing the objects like mirrors, lenses, prisms, aperture and
even an opaque body in the path of the light.
The details of the image can be obtained either by applying
1) the principles of geometrical optics in which a beam of rays of light is considered
(or)
2) the principles of physical optics in which wave nature of light is considered.
The selection of the particular branch among these two depends on three factors namely
1) the wavelength of the light - ()
2) the size of the object - (b)
3) the distance of the screen from the object - (l)
The condition that should be satisfied for the use of geometrical optics is b >>
The condition that should be satisfied for the use of physical optics is b » (or) b <
Huygen’s Wave Theory:
From the Huygen’s wave theory it follows that the light energy is propagated in the form of
waves. Fresnel suggested that these waves are transverse in nature. A point source of light is supposed to emit
a wave which spreads equally in all directions in an isotropic medium. Such a disturbance is called a spherical
wave.
Wave front: The continuous locus of all the points at which the vibrations are ‘isophasal’ is called the
wave front. The word isophasal indicates that the phase difference between any two points on the wave front is
zero (and is not 2p or 4p). However the phase difference between the particles of one wave front and the
particles of another wave front at a distance ‘‘ apart is 2. The phase difference between the two wave fronts
x
at a distance x apart is 2. .
Huygen’s Principle: According to Huygen’s principle every point on a wave front acts as a secondary
source of disturbance and sends out wavelets in the forward direction with a velocity equal to the velocity of light.
In the diagram, it is shown how the wavelets are derived from a slit obstructing the wave front.
Interference:
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IIT - FOUNDATION - SET - V
The light from a single source is distributed uniformly in the space surrounding it. But if there are two
sources then the waves from one source “interfere” with the waves from the other source and the distribution of
light becomes non-uniform. The intensity at any point depends on how the waves from these sources arrive at
that point.
At a point where the crests of the two waves (or) the troughs of the two waves always arrive simulta-
neously, the two waves reinforce. The intensity at such a point becomes larger. This is called constructive
interference.
At a point where the crests of one wave arrive exactly at the instant, when the troughs of the other wave
arrive then they destroy each other. The intensity at such a point becomes zero. This is called destructive
interference.
Thus as a result of interference of one wave with the other, the distribution of light becomes non-uniform
and the space surrounding the source is divided into alternate bright and dark regions. Under certain conditions
this becomes a ‘permanent’ pattern.
Principle of superposition:
The net displacement at any point due to waves from a number of sources is given by the
algebraic sum of the displacements due to the individual waves.
If y1, y2, y3........... are the displacements at a given point respectively due to a number of waves
then the net displacement at that point is given by y = y1 + y2 + y3 +..............
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PHYSICS
Theory of interference:
Let S1 and S2 be two monochromatic sources of same frequency . The wavelength of the
wave from each source is . Let P be a point at distances x1 and x2 from the sources S1 and S2. The
displacement y1 at P due to the wave from the source S1 is given by
y1 = a sin(t - kx1) …………..(1)
where a = amplitude, = 2, kx1 = phase constant of the wave and k = 2/
Similarly, the displacement y2 at P due to the wave from S2 is given by
y2 = a sin(t - kx2) ………..(2)
From the principle of superposition, the net displacement ‘y’ is given by
y = y1 + y2
y = a[sin(t - kx1) + sin(t - kx2)]
kx kx1 kx2 kx1
y = 2a.cos 2 .sin t …………… (3)
2 2
kx2 kx1
The term 2a.cos represents the net amplitude at P. (kx2 - kx1) gives the phase
2
difference between the waves and (x2 - x1) gives the path difference between the waves.
The amplitude becomes a maximum of ‘2a’ and hence the intensity at ‘P’ becomes
kx2 kx1
maximum when cos =1
2
kx2 kx1
(or) = 0, , 2, 3, ............
2
k(x2 - x1)= 0, 2, 4, 6, ...... i.e. phase difference = 2n, where n = 0,
1, 2, 3, ......
2
(or) (x - x ) = 2n
2 1
(or) (x2 - x1) = 0, , 2, 3, ............
i.e. path difference = n, where n = 0, 1, 2, .....
This is the condition for constructive interference.
The intensity is proportional to the square of the amplitude. The maximum value of the intensity
of the resultant wave is Imax = (2a)2 = 4a2 = 4I.
when n = 0, path difference = 0. This corresponds to the central bright fringe.
kx2 kx1
Similarly, the intensity at P becomes minimum or zero when cos =0
2
kx2 kx1 3 5
i.e. = , , ,......
2 2 2 2
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IIT - FOUNDATION - SET - V
k(x2 - x1) = , 3, 5, ...........
i.e. phase difference = (2n - 1) , where n = 1, 2, .........
2
x2 x1 = (2n - 1)
2
(x2 - x1) = (2n - 1)
2
i.e. path difference = (2n - 1) , where n = 1, 2, .........
2
This is the condition for the destructive interference.
Thus, the condition for the constructive interference is Path difference = nl
(or) Phase difference= 2n.
The condition for the destructive interference is Path difference = (2n - 1) /2
(or) Phase difference= (2n-1).
Also, the two sources should have the same frequency and the two waves should have a constant phase
difference for the interference pattern to remain permanent. Two such sources satisfying these conditions are
called the coherent sources.
Coherent Sources :
Definition: Two monochromatic sources are said to be coherent if they have the same frequency and
bear a constant phase difference.
1) Similarly two waves are said to be coherent if they have same wavelength and bear a constant phase
difference.
2) Two different sources can never be coherent, though they have the same frequency because there is
no guarantee that the phase difference between them remains constant. Similarly, waves derived from two
different points of an extended source cannot be coherent. From the Huygen’s wave theory a wave front is an
isophasal surface. Every point on the wave front acts as a secondary source and sends out wavelets in all forward
directions. So a wave front from a single point on a monochromatic source should be divided into two parts and
these two parts can be taken to be coherent.
Definition of Interference: Interference is that phenomenon in which the waves arriving from two co-
herent sources superpose one over the other in certain region and produce a permanent interference pattern
consisting of alternate bright and dark fringes in that region.
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PHYSICS
Young’s double slit experiment:
A) Description: A narrow slit ‘S’ is illuminated by using a monochromatic source of light M. Two more
narrow slits S1 and S2 having same width are held parallel to each other so that they are situated symmetrically
with respect to slit ‘S’. These slits S1 and S2 should be at a small distance ‘d’ apart. A screen is held at a distance
‘D’ from these slits.
B) Phenomenon Observed : A number of alternate bright and dark fringes are noticed. All these
fringes are found to be of equal width and are all parallel to each other. The central fringe happens to be a bright
fringe.
C) Theory : The cylindrical wave front from the slit S is incident on the slits S1 and S2. The wavelets that
pass through the slits S1 and S2 are derived from the same wave front. So they are coherent.
Let P be a point on the screen which is at distances x1 and x2 from the coherent sources and . Let
y1 and y2 be the displacements at P due to the wavelets arriving from S1 and S2.
y1 a.sin wt kx1
2
y2 a.sin wt kx2 where k and ‘’ is the wavelength of the light.
k x2 x1
The net displacement ‘y’ at the point P is given by y 2a.cos .sin
2
k x2 x1
wt
2
k x2 x1
Amplitude of the resultant wave at ‘P’ = 2a. cos
2
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IIT - FOUNDATION - SET - V
But K(x2-x1) is the phase difference between the wavelets arriving at P.
and (x2-x1) is the path difference between the wavelets arriving at P.
The square of this amplitude represents the intensity of the light at P.
The square of this amplitude represents the intensity of the light at P.
1) Condition for having bright fringe at P.
The intensity at P will be maximum and equal to 4a2 when
K x2 x1
Cos2 =1
2
or K(x2 -x1) = = 0 or 2 or 4 .......
or Phase difference = n.2 where n = 0, 1, 2, .....
or Path difference (x2-x1) = n where x = 0, 1, 2 ......
2) Condition having dark fringe at P.
The intensity at P will be minimum and equal to zero when
K x2 x1
cos2 =0
2
2
or K (x2-x1) = ( x2 x1 ) (x2-x1) = or 3 or 5.....
or Phase difference = (2n - 1) /2 where n = 1, 2, 3, ......
or Path difference (x2 - x1) = (2n-1) /2 where x = 1, 2, 3, ......
Thus alternate bright and dark fringes appear on the screen. For these fringes to appear permanently at
their respective positions, the amplitude must be independent of time. That is why the two sources or the two
waves are required to be coherent bearing a constant phase difference.
3) Fringe width: The distance between two consecutive bright fringes or the distance between the
two consecutive dark fringes is called the fringe width. It can be shown that fringe width
D
= where D = distance of the screen from the two slits and d = distance between the
d
two slits.
4) In general if ‘a1’and ‘a2’ are the amplitudes of the interfering waves then their intensities can be
taken as I1=a12 and I2 =a22.
The amplitude of the resultant wave at P is given by a2 = a12 + a22 +2a1a2.cos where is the
phase difference between the waves and =k (x2 – x1)
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PHYSICS
The intensity at P is therefore given by I = I1 + I2 + 2 cos.
When = 0 then a2 = (a1 + a2 )2 and I = maximum.
When = then a2 = (a1 a2 )2 and I = = minimum.
Applications of Interference :
1) A thin transparent film placed in the path of one of the interfering waves introduces an additionalpath
difference, which depends on the thickness of the film and the refractive index of the material of the
film. A method is developed based on this property which helps us to find accurately (a) thickness of
the film and (b) the refractive index of the material of the film.
2) A similar method is also evolved to find the velocity of light in different gases maintained at different
pressures.
3) A method involving interference is also developed to check whether a given surface is optically flat or
not.
4) Interference of light is used to determine the wavelength of sodium light and the wave length differ-
ence between D1, D2 lines.
5) Refractive index of liquids and gases can be very accurately determined by interference technique.
6) The reflecting power of the lens and prism surfaces can be tested by means of interference.
7) Michelson interferometer is used to determine velocity of light.
Diffraction:
Huygen’s wave theory gives explanation for the rectilinear propagation of light. However, a deviation
from rectilinear propagation was noticed in Grimaldi’s experiment. He formed the shadow of an obstacle on a
screen and noticed that (a) some light has encroached into the geometric shadow and (b) the edge of the shadow
is not sharp and surrounded by a number of alternate bright and less bright fringes. This opens a new topic called
diffraction of light.
Similar examples of diffraction: (1) the silver lining of a cloud (2) the silver lining around the profile of
a hill before sunrise. (3) Beautiful diffraction patterns of light seen through overlapping screens with tiny perfora-
tions and the change of such patterns while one perforated screen is moving or rotating relative to the other.
Definition of diffraction: The encroachment of light into geometric shadow of an obstacle and the
variation in the intensity of light around the edge of the shadow constitute a new phenomenon called diffraction.
Fresnel gave explanation to diffraction phenomenon treating it as the result of the mutual interfer-
ence of the secondary wavelets, which started from the different zones of the wave front as it is obstructed.
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Types of Diffraction:
The diffraction phenomena are classified into two heads.
1) Fresnel class of diffraction and (2) Fraunhofer class of diffraction.
If the wavelengths of the light used is ‘ ’ , the size of the object is ‘b’ and the distance (real or virtual)
between the obstacle and the screen is ‘’ then
1) the condition to classify the phenomena as Fresnel’s diffraction is thatb
2) the condition to classify the phenomena as Fraunhoffer’s diffraction is that b < <
Fresnel’s diffraction at a straight edge :
(a) Arrangement: A narrow slit S is illuminated by a source of monochromatic radiation. The cylindrical
wave front AB is obstructed by a plate EF with a straight and sharp edge E. The point ‘O’ on the screen is in line
with SE.
(b) Phenomenon noticed on the screen: We expect a geometric shadow below ‘O’ on the screen
and general illumination above ‘O’. But contrary to our expectation bright and less bright bands of unequal width
are noticed above ‘O’. And immediately below ‘O’ in the geometric shadow some light is seen encroached into
the shadow region the brightness of which decreases quickly and thereafter only darkness exists.
(c) Fresnel’s Explanation: The cylindrical wave front AB that is generated from the slit is obstructed
by EF.
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PHYSICS
Intensity at point ‘P’ above ‘O’ : Fresnel suggested that the wave front can be divided into two halves
AM and MB with respect to the point P. This point M is called the pole with respect to the point P. He further
suggested that each half could be divided into zones. MC1, C1C2 C2C3..........are the zones in the upper half and
MD1, D1D2, D2D3...... are the zones in the lower half.
This division is so made that the wavelets arriving at P from any two successive zones differ in phase by
‘’. So these zones are called the half period zones. He further made the following assumption regarding these
wavelets.
As one gets away from the pole ‘M’ the amplitude of the wavelets that arive at P from the successive half-
period zones decreases gradually.
Based on this assumptions Fresnel showed that the intensity at the point P in the illuminated region on the
screen (a) becomes maximum if the region ME of the lower half has odd number of half period zones (b)
become minimum if there are even number of half period zones in this region. (This minimum is not zero).
Thus in the region above ‘O’ fringes of unequal width which are alternately more bright and less bright are
formed. He further showed that the difference in the brightness of these fringes decreases very quickly with the
increase of the distance of the point ‘P’ from ‘O’. So only few such fringes are noticed and thereafter there is only
general illumination.
Illumination at a point ‘Q’ below ‘O’ that is in the shadow region :
‘N’ is the pole w.r.t the point ‘Q’. The lower half NB is completely cut off by the obstacle. In the upper
half the portion EN is obstructed. So the point ‘Q’ is exposed to the portion AE only in the upper half. From the
Fresnel’s assumptions it can be concluded that some illumination can be noticed at a point ‘Q’ in the
shadow region and this illumination decreases very quickly with the increase of the distance of ‘Q’ from ‘O’. So
in a small distance away from ‘O’ the illumination becomes almost zero.
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Fraunhofer diffraction due to a single slit :
The monochromatic light from the slit ‘s is rendered parallel using the lens L1. This parallel beam of light
otherwise called the plane wave front is intercepted by a slit AB of a very small width. The emergent beam is
brought to focus, using another lens L2, on a screen. The wavelets that are derived from different points of the
portion AB of incident wave front undergo constructive interfernece and give rise to central maximum of large
intensity at ‘O’ on the screen.
But at a point like ‘P’ on either side of ‘O’ the wavelets will either undergo constructive or destructive
interference depending on the value of ‘q. Thus on either side of the central maximum at ‘O’ a number of alternate
maxima and minima are formed. The intensity at any such minimum is zero. The intensity at any secondary
maximum is very small compared to the intensity of central maximum.
Applications of diffraction:
1. The study of Fraunhofer diffraction through a slit helped to make an optical component called the
‘diffraction grating’ with the help of which the spectrum of a composite light can be formed and the
wavelengths of the different spectral lines can be determined.
2. The study of the diffraction of X-rays through the spacing between the different planes containing the
ions of a crystal helps us to understand the (a) structure of the crystal and (b) to find the wavelength
of the X-rays.
3. Methods involving the diffraction of ultrasonic sounds are developed to find the velocity of sound
through liquids.
Comparison of Fresnel and Fraunhoffer diffraction :
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PHYSICS
Fresnel’s Diffraction Fraunhofer’s Diffraction
1. The wave front is curved i.e. 1. The wave front is plane.
cylindrical or spherical.
2. Source and screen are at finite 2. The effective distances of the source
distance from the obstacle. and the screen are infinity.
3. It is position dependent. 3. It is direction dependent.
4. No lenses are used for 4. Lenses are needed for observation.
observation.
5. It can be seen with naked eye. 5. Telescope focused to infinity is
necessary to observe the fringes.
6. Fringes are faint. 6. Fringes are bright.
7. b 7. b <<
Polarisation:
Unpolarised and Polarised light:
According to Maxwell’s electromagnetic theory, light consists of simple harmonic disturbances
in the strengths of the electric and magnetic fields which are at right angles to each other and at right angles to the
direction of propagation of light. In the case of ordinary light, the directions of electric and magnetic field vectors
at any point, orient at random. In other words the vibrations take place in all directions in the plane perpendicular
to the direction of propagation.
In the case of polarised light, the vibrations take place in a fixed plane and at right angles to the direction of
propagation. This is the reason why the polarisation is said to be simply - one sidedness.
The plane in which the vibrations in the intensity ‘E’ of the electric field are taking place is called the plane
of vibration and the plane perpendicular to it is called the plane of polarisation.
The rays of unpolarised and polarised light are represented as follows:
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Polarisation of light accompanied by reflection and by refraction:
Malus found that during the reflection of light at a surface backed by a transparent medium, the reflected
light in general is partially polarised. However the reflected light will be completely polarised for certain angle of
incidence, called the angle of polarisation. The refracted component at the same time is not completely polarised
and is only partially polarised.
Ex : For glass, the angle of polarisation is 570.
For pure water, it is 530.
Also it is found that when light is incident at an angle equal to the angle of polarisation on a surface
backed by a transparent medium, both reflection and refraction take place. The reflected light is polarised with
vibrations at right angles to the plane of reflection. The refracted light has a polarised component with vibrations
in the plane of refraction. Also, the reflected and refracted rays are found to be at right angles.
Brewster’s law: If the angle of polarisation is for the surface of a medium of refractive index
then tan =
The angle of polarisation for a given material depends on the colour of the light.
Double Refraction:
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Bertholinus discovered that when a ray of light is incident on a crystal like calcite, there are two refracted
rays and found further that -
One of these refracted rays obeys the Snell’s law of refraction and has a fixed refractive
index 0. This ray is called ordinary ray or O-ray. The value of 0 is independent of direction.
The other ray does not obey the Snell’s law and has no fixed value for its refractive index
e. This ray is called extraordinary ray or E-ray. The value of e depends on the direction.
For crystals like calcite otherwise called negative crystals, e 0
For crystals like Quartz otherwise called positive crystals, e 0
The ordinary and extraordinary components are both plane polarised, with their vibrations
at right angles to each other.
The wave front of O-rays is spherical and the wave front of E-rays is elliptical. So velocity
of O-ray is same in all directions, while the velocity of E-rays is different in different directions.
The crystal is isotropic for O-ray and anisotropic for E-rays.
In Calcite Ve V0, In quartz Ve V0
However, both wave fronts travel with the same velocity in certain direction, called
the optic axis of the crystal and along this direction e = 0
Production of Polarised Light:
Method (1) : By reflection and refraction using pile of plates: When light is incident at Brewster’s
angle on a plate of glass then
a) The reflected component is entirely plane polarised with vibrations perpendicular to the plane of
reflection but its intensity is only 7% of the incident light.
b) The refracted component consists of two parts. Part-1 has vibrations perpendicular and part-2 has its
vibrations parallel to the plane of refraction.
So, part-1 of the refracted component should also be made reflected so that the reflected component
becomes more intense and the refracted component consists of only part-2. This makes the refracted compo-
nent completely plane polarised. This is achieved by arranging a pile of parallel glass plates in the path of the
incident light i.e., by making the light to undergo successive reflections at all these plates. It is estimated that a pile
of 32 plates will give good result.
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Method (2): By double refraction :
1) Dichroic crystals: The dichroic crystals have the property of selectively absorbing one of the beams
produced by double refraction. This property is called dichroism. For example, Tourmaline crystal of thickness
2mm absorbs the O-rays completely. So the emergent beam from such a crystal consists of E-rays which is
completely plane - polarised.
2) Polaroid: A polaroid is a thin film of cellulose which are impregnated with minute dichroic crystals of
organic compounds, the optic axes of all the crystals being set parallel to each other. If unpolarised light is
incident on one face, polarised light emerges from the other face.
Note: If two polaroids or Tourmaline plates are arranged with their axes at right angles and light is made
to incident, no light is seen transmitted. This is because, the first polaroid absorbs the O-rays and transmits the
E-rays. The E-rays act as O-rays for the second polaroid. So these E-rays are also absorbed by second
polaroid. Therefore, no light is transmitted.
To distinguish practically the plane polarised light from unpolarised light:
The given light is seen through a polaroid. If the intensity remains the same on rotating the
polaroid then the given light is unpolarised. Because, if the given light is unpolarised then it has always one
component which acts as E-ray and is allowed to pass through the polaroid, in all its positions.
If the intensity changes and becomes alternately maximum and minimum on rotating the polaroid
then the given light is plane polarised. Because in some positions of the polaroid, the plane polarised light is
transmitted through it as E-rays. On rotating the polaroid through 900 the given polarised light becomes O-rays
for the polariod and hence is absorbed by the polaroid. Thus the intensity alternately becomes maximum and
minimum.
Uses of Polarised light:
1) Suppressing the glare from a polished surface: The glare due to the reflection from a polished
surface can be avoided by looking through a polaroid. The reflected light is partially polarised. The polarised
component is absorbed by the polaroid in certain position and the rest of the light only is allowed to pass through.
This avoids the glare.
2) Glareless head lights: The glare of the head light of the car coming in the opposite direction is the
major reason for the accidents during the night.
This can be avoided by covering the head light and windscreen glasses in front of the driver with
polaroids. For the cars approaching in opposite direction, the polarised light from the head lights of one car
become O-rays for the wind screen of the other car. So this light is absorbed by the wind screen and does not
reach the driver of the car. At the same time the driver of one car can see the other car under his own head light.
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PHYSICS
3) Stereoscopic or 3-D pictures: In 3-D pictures, the scene was shot using two cameras held at some
angle. So two projectors are used in the theatre. But the lenses of these projectors are covered with polaroids
in such a way that the two beams are polarised at right angles. If now the image on the screen is seen through
spectacles made of polaroids arranged at right angles, then the depth of the scene is felt. In other words, we have
a 3-D vision.
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ASSIGNMENT
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PHYSICS
1. The intensity of the central maxima in Young’s double slit experiment is I. If one of the slits is closed the
intensity at the same point will be
I I
1) 2) I 3) 0 4)
2 4
2. If white light is used in young’s experiment
1) all the fringes appear coloured
2) only bright fringes appear coloured
3) central fringe appears white and the remaining bright fringes appear coloured
4) None
3. If the intensities of the two sources are in the ratio 9:4, then the ratio of maximum to minimum intensity
ratio in the interference pattern is
1) 5:1 2) 25:1 3) 9:4 4) 81:16
4. In Young’s double slit experiment, the ratio of intensities of maxima and minima is 9:1, this implies that
1) the intensities at the screen due to two slits are 5 and 4 units
2) the intensities at the screen due to two slits are 3 and 2 units
3) their amplitude ratio is 3 : 2
4) their amplitude ratio is 2 : 1
5. Two beams of light having intensities I and 4I interfere to produce a fringe pattern on a screen. The
phase difference between the beams is /2 at a point A and p at point B. Then the difference between
the resultant intensities at A and B is
1) 2I 2) 4I 3) 5I 4) 7I
2
6. The path difference between two waves meeting at a point is . If both of them have intensity I, then
3
the intensity of the interference pattern at that point is
1) I 2) 2I 3) 3I 4) 4I
7. The intensity at the centre of the screen in ideal (coherent sources of same intensity) Young’s double slit
experiment is I0. The intensity at a point on the screen where the phase difference between the waves is
/2, is
I0 I0
1) 2I0 2) 3) 4) I0
4 2
8. In Young’s double slit experiment the two slits act as coherent sources of equal amplitude and of wave-
length l. In another experiment, the two slits act as incoherent sources. The ratio of intensity of light at
the centre of screen in the first case to that in the second case is
1) 1:1 2) 1:2 3) 2:1 4) 4:1
9. A glass plate of thickness 1mm is introduced in one of the interfering beams of young’s arrangement. If
0
m = 1.5 and l = 5000 A , The number of fringes shifted is
1) 500 2) 1000 3) 2000 4) 100
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o
10. A monochromatic light of wavelength 5000 A is incident on two slits separated by a distance of 5x10-
4
m. If a thin glass of thickness 1.5x10-6m and refractive index 1.5 is placed between one of the slits and
screen. The phase difference introduced at the position of the central maximum is
1) 3p 2) 2p 3) p/2 4) p/3
0
11. The fringe width in air in young’s experiment is 3.6mm for light of l=6000 A . If the arrangement is
immersed in a transparent liquid of =1.2 the fringe width would be
1) 4.32mm 2) 3mm3) 0.33mm 4) 1.8mm
12. The angular separation of two bright fringes in air is 0.240. If the arrangement is immersed in water
4
the angular separation becomes
3
1) 0.320 2) 0.180 3) 0.240 4) none
13. In Young’s double slit experiment, the distance between the slits is d. The screen is at a distance D from
slits. If a bright fringe is formed opposite to a slit on the screen, the order of the fringe is
d2 2d d 2d
1) 2) 3) 4)
2D d2 D D
14. In Young’s double slit experiment the 16th bright fringe of wavelength 600nm is formed at a certain
place. At the same place if 20th bright fringe is formed, the wavelength of the light used is
1) 240nm 2) 400nm 3) 480nm 4) 500nm
15. In Young’s double slit experiment the distance between the slits is 0.9mm. and the distance of the screen
is 1m from the slits. If the second dark fringe is at a distance of 1mm from the central maximum then the
wave length of the light used is
1) 400nm 2) 500nm 3) 600nm 4) 700nm
16. In two separate Young’s double slit experiments the ratio of the slits separation is 2 : 1 and the ratio of
the wavelengths used is 1 : 2. The fringe width is same. Then the screen distances from the slits in the
two experiments are in the ratio
1) 4 : 1 2) 1 : 4 3) 1 : 1 4) 2 : 1
17. Two glass plates of refractive indices 1.5 and 1.6 are introduced in the paths of two interfering beams in
Young’s experiment. The central fringe is not shifted. The thickness of the plates are in the ratio
1) 6 : 5 2) 5 : 6 3) 6: 5 4) 16 : 15
18. Light of wavelength l in air enters a medium of refractive index m. The phase difference between two
points in this medium, lying along the path of this light at a distance x apart is
2 1 2 2 1 2
1) x 2) x 3) ( 1) x 4) x
( 1)
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19. Shadows of aeroplanes and birds flying high are not formed on the surface of the earth. This is due to
1) interference 2) diffraction of light
3) polarisation 4) refraction
20. The class of diffraction in which the incident and diffracted wavefronts are either spherical or cylindrical
is called
1) Fresnels diffraction 2) Fraunhofer diffraction
3) Huygens diffraction 4) None
21. In Fresnel’s zones the average phase difference between secondary waves from two consecutive zones
reaching a point on the screen is
1) zero 2) /2 3) 4) 2
22. In Fresnel’s zones the operative zones contributing intensity are
1) last zones 2) first few zones 3) middle zones 4) all the zones
23. The average path difference between two waves coming from third and fifth fresnel zones of a wave
front at the centre of the screen is
1) 2) 2 3) l 4) 4l
2
24. A ray of light is incident on a glass plate at an angle 600. It was found that the reflected and refracted
rays are at right angles to each other. Then the refractive index of glass used is
1) 1.5 2) 1.414 3) 1.732 4) 2.0
25. A ray of light is incident on the surface of a glass plate at an angle of incidence equal to Brewster’s angle
f. If m represents the refractive index of glass with respect to air, the angle between the reflected and the
refracted rays is :
1) (90+f) 2) sin1(m cos f) 3) 900 4) sin 1(sinf/m)
26. An ink dot marked on a piece of paper is observed through a calcite crystal. The number of images
observed is
1) one only 2) two only
3) three only 4) sometimes one and sometimes two
27. When unpolarised light passes through a polaroid sheet, the beam that emerges from it is plane polarised.
This is due to selective
1) absorption of the O-Ray 2) absorption of the E-Ray
3) absorption of both E & O Ray 4) reflection of one of the rays
28. When a beam of plane Polarised light passes through a rotating Nicol Prism, the intensity of the transmit-
ted light
1) is always zero 2) Passes through maximum and zero
3) remains same 4) Passes through maximum and minimum
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29. When a Beam of light is observed through a Nicol which is slowly rotated the intensity of light Passes
through maximum and minimum but not zero. The beam is
1) plane polarised 2) unpolarised
3) partially polarised 4) circularly polarised
30. An analysing nicol examines two adjacent plane polarised beams A and B whose planes of polarization
are mutually perpendicular. Beam B shows zero intensity. From this position a rotation of 300 shows
that the two beams have same intensity. The ratio of intensities IA : IB will be
1) 1 : 3 2) 3 : 1 3) 3 :1 4) 1 : 3
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PHYSICS
KEY
1 4 16 1
2 3 17 1
3 2 18 1
4 4 19 2
5 2 20 1
6 1 21 3
7 3 22 2
8 3 23 3
9 2 24 3
10 1 25 3
11 2 26 4
12 2 27 1
13 1 28 2
14 3 29 3
15 3 30 2
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CHEMISTRY
Contents
1. ATOMIC STRUCTURE
2. CHEMICAL BONDING
3. ACID AND BASES
4. SALTS
5. ELECTRSO CHEMISTRY
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1 ATOMIC STRUCTURE
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CHEMISTRY
Introduction
An atom consists of subatomic particles like electrons, protons and neutrons. Protons and neutrons remain
together at the centre of the atom in the form of nucleus and the electrons revolve around the nucleus. Since in an
atom the number of protons is equal to the number of electrons, the atom is electrically neutral.
According to Bohr’s atomic model, the electrons revolve around the nucleus in certain fixed orbits and
while doing so, they do not lose any energy. Bohr’s theory successfully established the stability of an atom.
During the discovery of the structure of the atom, an experiment was carried out by passing the energy
emitted by the excited hydrogen atom through a prism. It was found that the atomic spectra of hydrogen atom
consisted of discrete spectral lines unlike white light which gives a continuous spectra. The presence of discrete
spectral lines indicated that the electrons emit fixed quantity of energy. The experiment proved the particle nature
of electron.
Particle nature of electron
In 1913, the particle nature of the electron was established by Bohr in his atomic model with the help of
Planck’s quantum theory of radiation. He established that energy emitted or absorbed by the electron is quan-
tized. According to the theory, the energy emitted or absorbed is a whole numbered multiple of hn, where h is
Planck’s constant and n is the frequency of the radiation. The Bohr’s theory could successfully explain the
discrete line spectra of hydrogen atoms. Bohr introduced principal quantum number which represents the main
energy levels around the nucleus.
Sommerfeld’s elliptical model
Sommerfeld extended the Bohr’s theory so that the theory could explain the splitting of spectral lines of
hydrogen atom. Sommerfeld introduced the concept of elliptical orbits in addition to the concept of circular orbit.
Azimuthal quantum number and energy sublevels or subshells were introduced to explain this model.
Explanation of Zeeman and Stark effect and magnetic quantum number
The explanation for splitting of a single spectral line of hydrogen atom into a number of closely spaced lines
in the presence of a magnetic field or an electric field was given by Lande. He postulated that in the presence of
an external electric field or a magnetic field, an electron can have different space orientations for a particular
angular momentum. The magnetic quantum number was introduced to specify the further subdivisions ofsubshells
or energy sublevels.
Spin quantum number
When each line of the spectral series of hydrogen atom was observed through an instrument with high
resolving power, it was found that each line was actually a combination of a pair of lines.
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To account for this, Uhlenbeck and Goud-smit suggested that each electron rotates on its own axis either
in a clockwise or in an anticlockwise direction while revolving around the nucleus. The spin quantum number was
introduced to explain this observation.
Wave nature of electron
de Broglie in 1924, pointed out that, the electron behaves both as a particle as well as a wave, just like
light.
de Broglie derived an expression for calculating wave length (l) of the wave which is associated with the
moving electrons.
h
=
momentum of electron , where h is Planck’s constant.
h
Or =
mv
This equation is called de Broglie’s equation, this equation was experimentally verified by Davison and
Germer’s experiment in 1927. Hence the dual nature of electron, i.e., particle nature as well as wave nature was
established experimentally.
Atomic orbitals
Bohr’s postulates state that electrons move around the nucleus with a fixed velocity in well defined orbits.
But on the basis of wave nature of electron and Heisenberg’s uncertainty principle, Bohr’s model of atom was
found to be incorrect as it is impossible to know simultaneously the exact position and velocity of a subatomic
particle like electron. During this period, Schrodinger established the fundamental wave equation of a small
particle in terms of wave motion.
On the basis of this wave equation, the probability of finding an electron having definite amount of energy
in a given space or region around the nucleus became possible. Consequently, the concept of atomic orbital came
into existence.
Atomic orbital is a three dimensional region in space around the nucleus, where the probability of finding an
electron with specific energy is high. There are four types of orbitals – s, p, d and f
Types of orbitals present in K, L, M and N shells.
K shell (n = 1) only s orbital
L shell (n = 2) s and p orbital
M shell ( n = 3) s, p, and d orbitals
N shell (n = 4) s, p, d and f orbitals
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Representation of an orbital
Since atomic orbital is a three dimensional space which represents the probability of finding an electron, it
is difficult to give a pictorial representation of the orbital.
It is generally represented by a shaded figure where the intensity of shading is proportional to the probabil-
ity of finding an electron in that area.
Shape of the orbital
An orbital can be represented by the surface where there is maximum probability of finding an electron.
s-orbitals: These orbitals have only one orientation and they are spherically symmetrical. The spherical
surfaces where the probability of finding an electron is zero are called nodal surfaces or nodes. The number of
nodes increases as n increases. For an s-orbital, the number of nodes is given by (n – 1), where n indicates the
main energy level.
Z
Y
The shape of 1s orbital
p-orbitals: These orbitals are not spherical like s – orbitals. They have two lobes separated by a node at
the nucleus. There are three orientations for p-orbitals represented as px, py and pz. The subscripts x, y and z refer
to the co-ordinate axes along which the density of the electrons of the respective orbitals is maximum. The p-
orbitals are dumb-bell shaped.
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z z z
y y y
x x x
The boundary surface diagrams of the p orbitals
d-orbitals: There are 5 d-orbitals represented as dxy, dyz, dzx, d x 2 y2 , d z 2 . The orbital dxy, dyz and dzx
have four lobes separated by a node at the nucleus. The lobes lie symmetrically between the co-ordinate axes,
for example, the lobes of dxy lie symmetrically in between x and y axes.
The orbital d ( x 2 y2 ) also has four lobes along x and y axes separated at the nucleus. The orbital d z2 has a
unique shape with two lobes along z-axis and a ‘belt’ like space centred in the x-y plane.
z z z
z z
y y y
y y
x x
x x x
dyz dxz dxy
The boundary surface diagrams of five d-orbitals.
Rules for filling up orbitals
The study of the atomic spectra of an atom showed different spectral lines, this provided the information
about the relative energies of the atomic orbitals. The increasing order of the energy of different orbitals belonging
to different main energy levels was given by Moeller.
Aufbau principle
“Aufbau” is a German expression which means building up or construction. According to this principle,
“the orbitals are filled up with electrons in the order of their increasing energy”.
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CHEMISTRY
1S
2S 2P
3S
3P
3d
4S 4P
4d 4f
5S 5P
5d
5f
6S 6P
6d
7S 7P
8S
Moeller diagram
Pauli’s exclusion principle states that “it is impossible for two electrons in a given atom to have all the
quantum numbers identical”.
According to Hund’s rule of maximum multiplicity, pairing of electrons in a degenerate orbital starts
when all the subrobitals are filled with unpaired electrons. Hence pairing of electrons p, d and f orbitals start with
the introduction of 4th, 6th and 8th electrons.
Given below are certain rules for filling up of electrons in the orbitals.
(i) The maximum number of electrons in a main energy level is equal to 2n2, where n is the number
of the main energy level.
(ii) The maximum number of electrons in a subshell or orbital like s, p, d and f is equal to 2(2–! + 1),
where the value of –! is 0, 1, 2 and 3 respectively for s, p, d and f orbitals. Hence the maximum
numbers of electrons in s, p, d and f are 2, 6, 10 and 14 respectively.
(iii) As a working rule, a new electron enters that orbital where (n + –!) is minimum, for example, the (n
+ –!) value of 4s orbital is 4 + 0 = 4 and that for the 3d orbital is 3 + 2 = 5. Hence the electron enters
4s orbital before entering into 3d. In case (n + –!) value is equal for two orbitals, then the electron
enters that orbital whose n value is less. For example (n + –!) value of 3d = 3 + 2 = 5 and 4p = 4 +
1 = 5, in this case, the electron prefers to enter 3d orbital.
(iv) Each orbital (s, px, py, pz, dxy dyz etc.,) can accommodate a maximum of two electrons with opposite
spin according to Pauli’s exclusion principle.
The electronic configuration of an atom is written in terms of nlx notation where n indicates the main
energy level, –! indicates the subshells or orbitals
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(s, p, d, f) and x indicates the total number of electrons present in the subshell.
(v) Electrons enter the orbitals px, py, pz, dxy dyz etc in such a way that maximum number of unpaired
electrons remain in a subshell and the unpaired electrons spin in the same direction.
Hence in the subshell like p, d and f electrons prefer to occupy separate orbitals rather than getting
paired in an orbital. These orbitals are identical in energy in the absence of any electric or magnetic
field and are called as degenerate orbitals.
As a result of this, pairing starts taking place only with the entry of 2nd electron in s – orbital, 4th
electron in p – orbital, 6th electron in d-orbital and 8th in f-orbital.
(vi) Half filled or completely filled orbitals are more stable than any other orbital.
Hence, actual electronic configuration of some transition elements are different from the expected
electronic configuration.
Example:
Expected electronic configurations of chromium and copper are 3d44s2 and 3d94s2 respectively.
But their actual electronic configurations are chromium: 3d54s1 and
copper: 3d104s1
Electronic configuration of first to fifth period elements are given below.
P Grou
At. K
e p Eleme L M N O
No n =
r numb nt n=2 n=3 n=4 n=5
. 1
i er
o 2 3 4 4 5 5 5
d 1s 2p 3p 3d 4p 4d 5s
s s s f p s f
F
i
r IA H 1 1
s
t
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VIIIA He 2 2
IA Li 3 2 1
S IIA Be 4 2 2
e IIIA B 5 2 2 1
c IVA C 6 2 2 2
o VA N 7 2 2 3
n VIA O 8 2 2 4
d VIIA F 9 2 2 5
VIIIA Ne 10 2 2 6
IA Na 11 2 2 6 1
IIA Mg 12 2 2 6 2
T
IIIA Al 13 2 2 6 2 1
h
IVA Si 14 2 2 6 2 2
i
VA P 15 2 2 6 2 3
r
d VIA S 16 2 2 6 2 4
VIIA Cl 17 2 2 6 2 5
VIIIA Ar 18 2 2 6 2 6
IA K 19 2 2 6 2 6 1
IIA Ca 20 2 2 6 2 6 2
IIIB Sc 21 2 2 6 2 6 1 2
IVB Ti 22 2 2 6 2 6 2 2
VB V 23 2 2 6 2 6 3 2
VIB Cr 24 2 2 6 2 6 5 1
F VIIB Mn 25 2 2 6 2 6 5 2
o
VIIIB Fe 26 2 2 6 2 6 6 2
u
VIIIB Co 27 2 2 6 2 6 7 2
r
t IB Cu 29 2 2 6 2 6 10 1
h IIB Zn 30 2 2 6 2 6 10 2
IIIA Ga 31 2 2 6 2 6 10 2 1
IVA Ge 32 2 2 6 2 6 10 2 2
VA As 33 2 2 6 2 6 10 2 3
VIA Se 34 2 2 6 2 6 10 2 4
VIIA Br 35 2 2 6 2 6 10 2 5
VIIIA Kr 36 2 2 6 2 6 10 2 6
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P Grou
At. K
e p Eleme L M N O
No n =
r numb nt n=2 n=3 n=4 n=5
. 1
i er
o 2 3 4 4 5 5 5
1s 2p 3p 3d 4p 4d 5s
d s s s f p s f
IA Rb 37 2 2 6 2 6 10 2 6 1
IIA Sr 38 2 2 6 2 6 10 2 6 2
IIIB Y 39 2 2 6 2 6 10 2 6 1 2
IVB Zr 40 2 2 6 2 6 10 2 6 2 2
VB Nb 41 2 2 6 2 6 10 2 6 4 1
VIB Mo 42 2 2 6 2 6 10 2 6 5 1
F VIIB Tc 43 2 2 6 2 6 10 2 6 5 1
i VIIIB Ru 44 2 2 6 2 6 10 2 6 7 1
f VIIIB Rh 45 2 2 6 2 6 10 2 6 8 1
t IB Ag 47 2 2 6 2 6 10 2 6 10 0
h IIB Cd 48 2 2 6 2 6 10 2 6 10 2
IIIA In 49 2 2 6 2 6 10 2 6 10 2 1
IVA Sn 50 2 2 6 2 6 10 2 6 10 2 2
VA Sb 51 2 2 6 2 6 10 2 6 10 2 3
VIA Te 52 2 2 6 2 6 10 2 6 10 2 4
VIIA I 53 2 2 6 2 6 10 2 6 10 2 5
VIIIA Xe 54 2 2 6 2 6 10 2 6 10 2 6
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2 CHEMICAL BONDNG
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Introduction
The atoms of all elements except noble gases have an inherent tendency to form chemical compounds to
decrease the energy of the system and thereby pass on to a more stable system. To attain this stable system, bond
formation takes place between the atoms. Different types of bond formation takes place depending on the nature
of participant atoms.
Ionic bond
An ionic bond is formed between two atoms, one of which loses one or more electrons and the other
which gains these electrons to attain the noble gas configuration. The atom losing the electrons froms a positively
changed ion (cation) and the atom gaining electrons forms a negatively changed ion (anion).
The cation and anion are held together by the electrostatic force of attraction.
Lattice energy of ionic compounds
Ionic compounds are formed due to the transfer an electron/electrons from one atom to the other. Hence
two oppositely charged ions are formed in the process and an electrostatic force of attraction is developed
between them.
In addition to the force of attraction between the individual pair, forces of attraction also exist among the
oppositely charged ions of different pairs which bring them close to each other. As a result of this, repulsive forces
are also generated between the similar ions of different pairs.
But the cation-cation distance and the anion-anion distance (measured from their centre) of different pairs
is greater than the cation-anion distance of different pairs.
Hence the force of attraction is greater than the repulsive forces and the ions arrange themselves in such a
way that their potential energy becomes minimum. The arrangement of ions where they possess minimumpoten-
tial energy gives the crystal structure of the ionic compound. The potential energy of the ions is assumed to be
zero when they are separated by an infinite distance.
The lattice energy of an ionic crystal is defined as the amount of energy released when the oppositely
charged ions are brought closer from infinite distance to form one mole of stable crystal.
Covalent bond
The covalent bond is formed, when two atoms attain their nearest noble gas structure by sharing one or
more electron pairs. Each of the atoms contributes equal number of electrons towards the bond formation.
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A covalent bond can be formed between two identical atoms as well as between the atoms of two different
elements.
Bond formation between two identical atoms
Let us take the formation of hydrogen molecule as an example.
When two hydrogen atoms are brought close together, two opposing forces start
acting
(i) Proton-proton and electron-electron repulsion.
(ii) Proton-electron attraction.
A system becomes stable when it possesses the lowest possible energy. Hence, a H2 molecule is formed
when the two H atoms position themselves in such a way that they have the minimum possible energy. The
distance where the energy of the atoms which took part in the bond formation is minimum, is called the bond
length.
The total energy of this system is a function of distance between hydrogen nuclei as shown in the following
graph.
Potential energy (KJ/mol.)
–458
0.074 (H – H bond length)
Inter molecular distance, (nm)
In the H2 molecule, the electrons reside primarily in the space between the two nuclei where they are
attracted simultaneously by the protons present in both the nuclei.
Bond formation between two H-atoms takes place at a distance between where the proton-electron
attraction just balances the electron-electron repulsion and proton-proton repulsion. This is the way bond forma-
tion takes place between two identical atoms.
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CHEMISTRY
Equal sharing of electrons takes place in this type of bonding. Hence the bond between two identical
atoms is called nonpolar convalent bond. No separation of charge takes place in the molecule.
Polar covalent bond
A polar covalent bond is a bond formed between two non-identical atoms. When the elements do not
differ much in their electronegativities, transfer of electrons is not possible between the atoms. Then sharing of
electrons takes place. Since the two atoms differ in their capacity to attract the shared electron pair, unequal
sharing of electron density results. As a result, polarity is developed in the bond. Though it is a covalent bond, a
slight ionic character is imparted to the bond due to the electronegativity difference between the bonded atoms.
Due to the unequal sharing of electrons, a slight charge separation takes place in this type of bonding.
Hence fractional positive and negative charges are developed in a polar covalent molecule and the molecules are
said to have dipole moment.
Comparative study of properties of ionic compounds, polar covalent compounds and non-polar covalent
compounds
Covalent
Properties Ionic
Polar Non-polar
Liquids or gases or
soft solids
Mostly
Physical state Crystalline solids Exception: Silicon
gaseous
carbide, SiO2
which are hard.
Intermediate
Low melting
Melting and boiling High melting point between ionic and
point and
points and boiling point non polar
boiling point
compounds
Insoluble in
Soluble in polar polar
Soluble in polar
solvents, insoluble solvents and
Solubility solvents and non
in non polar soluble in
polar solvents.
solvents. non polar
solvents.
Conduct electricity Generally less
in fused state or in conducting. But in
their aqueous Non
Electrical conductivity dilute aqueous
solutions conductors
solutions
conductivity is
more.
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Orbital overlapping during the formation of covalent bond
A covalent bond is formed when an overlap takes place between the orbital of one atom with the orbital of
another atom. Each of the overlapping orbitals contains an electron each.
Conditions required for orbital overlapping
a. The orbitals belonging to the valence shell take part in overlapping.
b. The electrons in the overlapping orbitals should have an opposite spin.
Different types of overlapping
(i) s-s overlapping: Example: H2 molecule. Each hydrogen atom has only one electron(1s1) which is avail-
able for bonding.
In the formation of H2 molecules, 1s orbital of one atom overlaps with that of the other forming a single
covalent bond.
H 1s1 1
H H
H 1s1 1
s-s overlap
(ii) s-p overlapping: Example: HF molecule. It is formed by the overlapping of 1s of hydrogen atom and
2pz of fluorine atom or by s-p overlapping.
H 1s1 1
F
F [He] 2s2 2p H
1 1 1
1 1 1 1
s-p overlap
Half filled orbital
(iii) p-p overlapping: Example: Cl2, F2, Br2 molecules etc. Cl2 molecule is formed by the overlapping of 3pz
orbitals of two chlorine atoms.
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1 1 1
Cl 3s2 3p5 1 1 1 1
Half filled orbital
1 1 1
Cl 3s2 3p2 1 1 1 1
Cl Cl
p-p overlap
In the above three examples (H2, HF, Cl2), two atoms, each having one unpaired electron come together
and these unpaired electrons overlap to form a single covalent bond.
If the atomic orbitals of the atoms possess more than one unpaired electrons, then more than one covalent
bond can be formed.
Strength of a covalent bond
The strength of a covalent bond is related to the extent to which the two combining atomic orbitals can
overlap.
The greater the overlap between the atomic orbitals, the greater is the strength of the resulting covalent
bond.
According to the above concept of the strength of the covalent bond, covalent bonds can be divided into
two types.
(i) Sigma bond due to end to end overlapping.
(ii) Pi bond due to sidewise or lateral overlapping.
Sigma () bond: Sigma bond is a strong covalent bond formed as a result of maximum overlapping (end to
end overlapping) of s-s, s-p and p-p orbitals along the internuclear axes.
Pi () bond: A weaker covalent bond which is formed between two atoms by lateral overlapping of
orbitals along a line perpendicular to the internuclear axis is called a pi bond.
Comparative study of sigma and pi bonds
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Sigma Bond Pi bond
1. It is a strong covalent formed due to 1. It is a weak covalent bond
end to end overlapping of orbitals. which is formed between two
atoms by lateral overlapping of
orbitals.
2. It is denoted by . 2. It is denoted by
3. It has an independent existence. 3. It does not have any
independent
existence as it can form only after
the formation of a bond.
4. Bond dissociation energy is high. 4. Bond dissociation energy is less.
5. Orbitals s, p and d are capable 5. Orbitals other than ‘s’ are
of forming bond. capable of forming a bond.
6. The bonded electron cloud is 6. The bonded electron cloud is
distributed in a cylindrically distributed in two banana shaped
symmetrical way around the regions above and below the
internuclear axis. internuclear axis.
+ Z Z
. .
Explanation of sigma and pi bond with the help of the formation of O2 and N2 molecules
O2 molecule: The electronic configuration of O-atom is 1s2, 2s2, 2px1, 2py2, 2pz1. Hence the oxygen atom
has 2 unpaired electrons. Singly occupied 2pz orbital of one oxygen atom will overlap (end to end) with the same
orbital of another Oxygen atom to give rise to a sigma bond. Singly occupied 2px orbitals of both the Oxygen
atoms overlap (side wise) along a line perpendicular to z-axis (molecular axis) to give a pi bond. Thus a double
bond is formed.
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px – px – bond
x x
y y
z ↿⇂ z
pz – pz – bond
O2 molecule ( bond – 1 and bond 1)
Thus an O2 molecule has a double bond between the atom, of the two bonds: one is a pz-pz sigma bond
and the other px-px pie bond called double bond.
N2 molecule: The electronic configuration of a nitrogen atom: 1s2, 2s2, 2px1, 2py1, 2pz1. Hence, a
Nitrogen atom has three singly occupied orbitals. One 2pz 2pz sigma bond is formed. The other
two pie bonds are 2px 2px and 2py 2py. Thus N2 molecule has three bonds: one pz pz sigma bond
and other two (p x p x and p y py) pie bonds which are perpendicular to each other and
perpendicular to the z-axis i.e. axis of the sigma bond. Hence, a triple bond is formed.
px – px –
xbond x py – py –
y
ybond
z ↿⇂ z
pz – pz – bond
N2 molecule ( bonds – 2 and bond – 1)
Geometrical shape of the covalent molecules
Geometrical shape of the covalent molecules is generated due to the repulsion between the
electron pairs present in the valence shell of the constituent atoms of the molecules.
Valence shell electron pairs present in the molecules arrange themselves in such a way that
repulsion between them becomes minimum.
In some molecules, non bonded valence shell electron pairs are present in the central atom.
These electron pairs are called lone pairs.
Repulsion between lone pair-lone pair is more than the repulsion between lone pair-bonded
pair and the repulsion between lone pair-bonded pair is more than the repulsion between
bonded pair-bonded pair. Hence the presence of a lone pair changes the shape of the molecules.
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Shapes of the molecules without a lone pair
Molecules Number of bonded pairs Shape
180°
BeCl2 2 Cl Be Cl (linear)
F
120°
120° B (triagonal
F planar)
BF3 3 120°
F
109.5°
C
(tetrahedral)
CH4 4 H H
H
Cl
90° Cl (triagonal
Cl P 120° bipyramidal)
PCl5 5 Cl
Cl
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Shapes of molecules having a lone pair
Number Number Change of shape due
Molecule Expected
of bond of lone to the presence of
s shape
pairs pairs a lone pair
Tetrahedra N
NH3 3 1
l H 107° H
H
(triagonal bipyramidal)
Tetrahedra P
PCl3 3 1
l 107° Cl
Cl
(triagonal bipyramidal)
Coordinate bond
In a coordinate covalent bond one of the atoms allows the other to provide both the electrons which are to
be shared. The other atom which accommodates the shared pair of electrons is called an acceptor. In this
mechanism, the contribution of electron pair is one-sided, but the sharing is equitable. Hence, slight polarity
develops in the molecule. A coordinate bond is represented as ‘’, pointing from the donor atom to the accep-
tor atom. A coordinate bond is explained below with the help of some examples.
(i) Formation of HClO:
.. .. .. .. H
..
H: .Cl. : .O. : H .Cl. .O. or H Cl .O. :
: : :
H
Cl-atom of HClO is the donor and O-atom is the acceptor of a lone pair of
electrons.
(ii) Formation of NH3 BF3
H F H F
H| F|
H N B F H N B F or H N B F
+ =
| |
H F H F H F
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N atom of ammonia has a lone pair of electrons and these are accepted by the vacant p orbital of the boron
atom of BF3. Hence N atom is the donor and boron is the acceptor.
Metallic bond
Each atom in a metal crystal loses all its valence-electrons. The electrons thus obtained form an electron-
pool. The resulting positively charged metal ions are believed to be held together by the electron pool. The
positively charged metal ions do not float randomly in the sea of electrons. They have a definite position in the
crystal lattice of the metal. The valence electrons are not attached to any individual ions. They belong to the
crystal as a whole and are free to move throughout the metal crystal.
The electrostatic force of attraction that binds metal ions to the mobile electrons within its sphere of
influence is known as a metallic bond.
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3 ACIDS & BASES
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THEORIES OF ACIDS & BASES:
Bronsted – Lowry acid - base theory:
ACID: Bronsted acid is a proton donor.
Ex : HCl H+ + Cl-
NH 4 H+ + NH3
H2 O
Here are Bronsted acids.
BASE: Bronsted base is a proton acceptor.
Ex : NH3 + H+ NH4+
H2O + H+ [H3O]+
Here are Bronsted bases.
NEUTRALISATION:
According to Bronsted-Lowry theory, neutralisation is the process in which a proton is transferred from an
acid to a base, forming conjugate acid- base pair
Ex 1 : HCl + NH3 NH4 + + Cl –
acid1 base2 acid 2 base1
Ex 2 : CH3COOH + H2O CH3COO– + H3O+
acid1 base2 base1 acid2
According to Bronsted - Lowry theory
Strong acids are those that have greater tendency to donate protons.
Ex : HCl , H2SO4 , HBr etc.
Weak acids are those that have fewer tendencies to donate protons.
Ex: CH3COOH, H2O , H3BO3etc.
A strong base is one that has greater tendency to accept a proton.
Ex : CH 3 COO , OH , CN
etc.
A weak base is one that has fewer tendencies to accept a proton and hold on to it.
Ex : NO3 . Cl , SO42 etc.
Conjugate acid - base pair:
Ans: An acid and base pair which differ by a proton is called a conjugate acid-base pair.
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acid – proton = base or base + proton = acid
Ex.(1) : HCl H+ + Cl (2) NH3 + H+ NH4+
acid c. base base c.acid
NH3 + H2O
base2 acid1 acid2 base1
NH4+ and NH3 are one conjugate acid base pair.
H2O and OH- are one conjugate acid base pair.
A weak acid has a strong conjugate base.
Ex.: CH3COOH CH3COO- + H+ ionised only partially.
weak acid strong conjugate base
A strong acid has a weak conjugate base.
Ex.: HCl H+ + Cl- (weak conjugate base).
(Strong acid)
Advantages and drawbacks of Bronsted-Lowry theory:
Advantages:
This theory explains the behaviour of acids and bases in both aqueous and non-aqueous solvents.
It is more generalised theory than Arrhenius which includes larger number of species in the list of acids
and bases.
Drawbacks:
A substance is said to be an acid in the presence of base only and viceversa.
It can not explain the acidic nature of electron deficient compounds like AlCl3, BCl3 etc.
Lewis-acid base theory:
ACCORDING TO LEWIS THEORY :
ACID : Electron pair acceptor is an acid
Ex : etc
The molecules (or) ions having vacant orbitals can accept an electron pair and
behave as lewis acids.
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CHEMISTRY
BASE : Electron pair donor is a base.
Ex : etc
The molecules or ions having lone pair of p electron pair can donate a pair of
electrons and behave as Lewi’s bases.
NEUTRALISATION:
According to Lewis , neutralisation is the process in which a coordinate covalent bond is formed between
a base and an acid
H3N H H 3 N H
Ex : H3N BF3 [H 3 N BF3 ]
Different types of Lewis acids:
Simple cations like Ag+, Al+3, Fe+2, Cu+2, Co+2 etc can accept the electron pairs and behave as Lewis
acids.
The molecules in which the central atom has incomplete octet (electron deficient compounds) are
also Lewis acids.
Ex : BCl3, BF3 , AlCl3 , FeCl3 etc
Molecules of compounds with central atoms having available d-orbitals to get more than an octet of
valence electrons
Ex : SiF4 , SF4 , SiCl4 , TeCl4.
Molecules with multiple bonds between atoms of dissimilar electronegativities
Ex : CO2 , SO2 , SO3 , NO2 , P4O10
Elements with electron sextet
Ex.: O, S.
Different types of Lewis bases:
The simple negative ions like Cl- , F- , OH- , CN- etc can donate the electron pair and
behave as Lewis bases.
Molecules with lone pairs of electrons
Ex: NH3 , H2O , R-O –R, R–O-H
Compounds with multiple bonds
Ex : H2C = CH2 , HC CH
Alkenes Alkynes
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Merits and demerits of Lewis theory:
Demerits:
i) This theory failed to explain the strength of acids and bases.
ii) It failed to explain the neutralisation of simple acids and bases, which do not involve in
Co-ordinate bond formation
Merits:
i) It increases the number of species that come under the category of acids and bases.
ii) It includes simple cations like Ag+ and Al+3 and simple anions like Cl-, OH- etc. in acids and bases
respectively.
Water acts as both Bronsted acid and Bronsted base:
Water molecule is capable of giving a proton and also accepting a proton. When HCl is dissolved in
water, H2O accepts a proton and acts as a base.
HCl + H2O H 3 O+ + Cl –
acid1 base2 acid2 base1
But when ammonia is dissolved in water, H2O donates a proton and acts like an acid.
NH3 + H2O NH 4 OH
Base2 acid1 acid2 base1
Thus water can act both as an acid and a base. Hence it is amphoteric.
Ammonia is a base according to both Lewis as well as Bronsted theory
In ammonia molecule, the central atom nitrogen has a lone pair of electrons. It can donate electron pair to
an acid and forms a co-ordinate covalent bond.
H3N: + BF3 [H3N BF3]
base acid
So ammonia is a base, according to Lewis theory. Also , Ammonia molecule can accept a proton to form
an acid.
NH3 + HCl NH4+ + Cl –
base1 acid2 acid1 base2
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So ammonia is a base, according to Bronsted theory.
All Bronsted – Lowry bases are Lewis bases
A bronsted base is one which accepts a proton.
Bromide ion (Br –) can accept a proton forming HBr molecule.
Br H
HBr Thus it acts as bronsted base.
A Lewis base is a molecule or ion having one or more lone pairs of electrons. So it can donate a pair of
electrons to an acid forming co-ordinate covalent bond.
..
: Br –: + AlBr3 [AlBr4] – Thus can also act as Lewis base.
Hence all Bronsted Lowry bases are Lewis bases.
Explain the following :
Most of the Lewis bases are Bronsted bases
A Lewis base is a molecule or ion having one or more lone pairs of electrons. So it can donate a lone pair
of electrons to an acid forming a co-ordinate covalent bond.
Ex : H3N : + BF3 [H3N BF3]. So, NH3 is a Lewis base.
Bronsted base is one which accepts a proton.
NH3 + H+ NH4+
So ammonia can act as Bronsted base also.
Thus, most of the Lewis bases are Bronsted bases.
All Lewis acids need not be Bronsted acids
Lewis acid is one, which has a vacant orbital and so can accept a pair of electrons forming a co-
ordinate covalent bond with a base
Ex :Cl3Al + Cl [Cl3 Al Cl] –
Bronsted acid is one which can donate a proton. In the above example the Lewis acid AlCl3 cannot give
a proton. So it cannot act as Bronsted acid. Therefore all Lewis acids need not be Bronsted acids.
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Q. What are the conjugate acids of the following?
(1) OH (2) NH 2 (3) HSO4 (4) HCO3
According to Bronste-Lowry theory Base + Proton = Conjugate acid.
Ion (base) Conjugate acid
OH H2O
NH 2 NH3
HSO 4 H2SO4
HCO 3 H2CO3
Q. What are the conjugate bases of the following?
(1) OH (2) NH3 (3) HSO4 (4) H2O
According to Bronste-Lowry theory Acid – Proton = Conjugate base.
Molecule or ion(acid) Conjugate base
OH O2-
NH3 NH 2
HSO4 SO42
H2 O OH
Levelling effect of water:
Water levels the strength of all strong acids to the strength of H3O+ and all strong bases to the strength of
OH– ions. It is called as levelling effect of water.
H3O+ is the strongest acid that exists in water and OH– is the strongest base that exists in water.
For example the acids HClO4, HCl, H2SO4 and HNO3 dissociate completely in water producing
H3O+ ions so all these acids show the same strength.
Q. a) Why Al+3 ion is a Lewis acid?
b) Pick out the Lewis base from the following?
FeCl3, AlCl3, BF3, NH3
a) Al+3 is a Lewis acid because, it has a vacant orbital to accept a pair of electrons from a base
to from a coordinate covalent bond.
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b) NH3 (because it has a lone pair of electrons on N) is Lewis base.
5. (a) What is neutralisation according to Lewis theory?
(b) Identify the Lewis acid from the following : NH3, AlCl3, H 2O, NaCl, Ag+.
(a) According to Lewis theory, neutralisation is the process in which a coordinate covalent bond is
formed between an acid and a base.
Ex: H3N: + H+ [H3N H ]+ or NH4+.
Cl 3 Al Cl Cl Al3 Cl or AlCl 4
acid
base
(b) AlCl3 , Ag+ (because of the presence of vacant orbitals.)
IONIC PRODUCT OF WATER:
Ionic product of water:
The product of hydrogen ion concentration and hydroxyl ion concentration in pure water (or) in aq.
solution at a given temperature is known as ionic product of water.
It is denoted by Kw.
OH 1.0 x 10
Kw H 14
mole 2 / lit 2
At 250C in pure water [H+] = [ OH ] = 1 x 10-7 moles/litre
At 250c, Kw = [H+] [ OH ]
Kw = 1 x 10-7 moles/litre x 1.0 x 10-7 moles /litre
= 1 x 10 -14 moles2 / litre2.
The value of Kw increases with increase in temperature.
pH and pOH:
pH of a solution:
pH of a solution is the negative logarithm to the base 10 of hydrogen ion concentration, expressed in moles
per litre in a given solution.
1
pH = -log10[H+] or pH log
H
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pOH of a solution:
pOH of a solution is the negative logarithm to the base 10 of hydroxyl ion concentration, expressed in
moles per litre in a given solution.
1
pOH = -log10[OH-] or pOH log
OH
P. What is the total value of pH and pOH in an aqueous solution at 250C.
pH + pOH = 14
Q. How does the pH value indicate the nature of an aqueous solution?
If the pH is less than 7, the solution is acidic and if the pH is more than 7, the solution is alkaline.
If the pH is exactly 7, then the solution is neutral.
In acidic solutions, with increase in pH the acidic nature decreases.
In alkaline solutions, with increase in pH the alkaline nature increases
Q. How do you predict the nature of an aqueous solution, from its H+ and OH– ion concentrations?
If the [H+] ion concentration is more than 110-7g. ions per litre, the solution is acidic.
Similarly if [H+] ion concentration is less than 1´10-7 g. ions per litre, the solution is basic.
If [ OH ] ion concentration is greater than 1´10-7g. ions/litre, and solution is alkaline.
If OH is less than 1´10-7g. ions/L. the solution is acidic.
Q. What is the importance of ionic product of water in aqueous solution? Classify the aqueous solutions
based upon the [H+] and [OH–].
From the value of ionic product of water, knowing either [H+] or [OH–], the other can be calculated.
In neutral solutions [H+] = [OH–]
In acidic solutions [H+] > [OH–]
In basic solutions [H+] < [OH–]
P. Find the pH of (a) 0.05 M.HCl solution (b) 0.01 M. H2SO4 solution.
+ -
(a) HCl
H + Cl , [ H ]= [HCl] since HCl is a strong acid.
[ H+] = 0.05 = 5´10-2 gram ions/litre.
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pH = -log [H+]
= -log [5 10-2] = 2 - log5 = 2-0.6990
pH = 1.3010
2
(b) H2SO4
2 H SO4
[H+] = 2 0.01
= 2 10-2 g. ions/litre
pH = –log 2 10-2 = 2 - log 2 = 2 - 0.3010 = 1.6990.
P. a) Calculate the pH of 0.01 M KOH solution
b) Find the pH of 0.001 M Ba(OH)2 solution.
Ans: a) [ OH ] = KOH
= 0.01 = 110-2; pOH = log OH = log 1 x 10 2 (2) 2
pOH = 2 and pH = 14 - 2 = 12.
b) Ba2+ + 2OH
Ba(OH)2
[ OH ] = 2 0.001 = 2 10-3 moles/litre.
pOH =
–log OH
= 3 - log 2
= 3 - 0.3010
= 2.6990
pH = 14 - pOH
= 14 - 2.6990
= 11.3010.
P. What is the hydrogen ion concentration of a solution whose (a) pH is 5, (b) pH is 9?
(a) I f the pH is 5, [ H +] = 10 pH or -anti log pH = -antilog 5
= 1 X 10-5moles/litre
(b) If the pH is 9, [H+] = 10-9 moles/litre
Q. a) What happens to the value of ionic product of water, when an electrolyte is dissolved in it ?
b) What is the nature of a solution whose pH is 6 if the ionic product of water at that temperature is
10-12 (mol /L.)2
a) The value of ionic product KW remains unchanged at a given temperature.
b) The solution is neutral. H OH 10 12 10 6 moles / litre
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P. Calculate the change in [H+] ion concentration if pH value of a solution decreases from 5 to 3.
Since [H+] = 10 pH
In 1st case when pH = 5, [H+] = 10-5moles/litre
In 2nd case when pH = 3, [H+] = 10-3 moles/litre
Change in [H+] is from 10-5 to 10-3 moles/litre
H+ ion conc. increases 100 times.
P. 50ml. of 0.2N H2SO4 is mixed with 100mL. of 0.2N HNO3 and diluted to 300mL.. What is the pH of
the solution?
V1 N 1 V 2 N 2 50 x0.2 100 x0.2
Normality of the mixture =
V 300
10 20
= 0. 1
300
pH = -log- [H+] = -log 0.1 = –log 10–1 = 1
P. The pH of an aqueous solution is 5.5. Calculate the [H+] ion concentration.
pH = -log [H+]
log [H+] = -pH = - 5.5 = - 6 +0.5 = 6 .5
Mantissa is always positive.
[H+] = Antilog of 6 .5
= 3.162 10-6 moles/litre.
BUFFER SOLUTIONS:
Q. What are buffer solutions ? Explain the buffer action of an acidic buffer in maintaining a steady pH
value of the solution ?
A solution whose pH remains unchanged even on dilution (or) on the addition of a small amount of a
strong acid (or) a strong base is called a buffer solution. (OR)
A solution of reserve acidity or alkalinity is called a buffer solution.
Acidic buffer is a mixture of solution of a weak acid and its salt with a strong base.
Ex : CH3COOH(aq.) + CH3COONa(aq.)
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Buffer Action : the above buffer solution is represented by the following equations.
CH 3COOH CH 3COO + H+
CH 3COONa CH 3COO Na
When a small amount of a strong acid like HCl is added, it ionises completely.
HCl H Cl
The added H+ ions are picked up by CH3COO- ions, present in large excess, forming unionised CH3COOH
molecules.
CH 3COO H CH 3COOH
Similarly when a small amount of a strong base like NaOH is added, it ionises completely.
NaOH Na OH
The added OH ions will react with CH3COOH molecules forming unionised water molecules and ac-
etate ions.
CH 3COOH OH CH 3 COO H 2 O
Thus the added H+ ions or OH- ions are removed from the solution. Therefore the PH of the solution
remains unchanged.
Q. What are buffer solutions ? Explain the buffer action of basic buffer in maintaining a steady PH value
of the solution ?
A solution whose PH remains unchanged on dilution or on the addition of a small amount of a strong
acid or a strong base is called a buffer solution.
A basic buffer is a mixture of solution of a weak base and its salt with a strong acid.
NH 4 OH NH 4 Cl
Ex :
aq.
aq.
Buffer action :
The above buffer solution is represented by the following equations,
NH 4 OH NH 4 OH
NH 4 Cl NH 4 Cl
When a small amount of a strong acid like HCl is added, it ionises completely.
HCl H Cl
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The added H+ ions react with NH4OH molecules forming unionised H2O molecules.
NH 4 OH H NH 4 H 2 O
Similarly, when a small amount of a strong base like NaOH is added, it ionises completely.
NaOH Na OH
The added OH ions are picked up by ions present in large excess, forming unionised NH4OH molecules
NH 4 OH NH 4OH
Thus the added H+ or OH- ions are removed from the solution. Therefore the pH of the solution remains
unchanged.
P. A buffer solution contains 0.01 moles/ litre of acetic acid and 0.1 moles/litre of sodium acetate. The
dissociation constant of acetic acid is 1.75 x 10-5. Calculate the pH of the buffer.
Henderson’s equation for pH of a acidic buffer is
[ salt ]
pH= P Ka + log
[acid ]
P K a of acetic acid = - log (1.75 x 10-5)
= 5 - 0.243
= 4.757
[salt] = 0.1 M ;
[acid] = 0.01 M
pH = 4.757 + log
= 4.757 + log 10
= 4.757 + 1
= 5.757.
P. What is the pH of the buffer solution obtained by mixing 20mL. of 0.2M CH3COOH and 10mL. of
0.4M sodium acetate? ( of CH3COOH = 4.8).
Given buffer is acidic.
0.4 x10
salt 30
pH = + log = 4.8 + log 20 x 0.2 = 4.8 + log = 4.8.
acid
30
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P. Calculate the pH of the buffer solution containing 20mL. of 0.1M NH4Cl and 2mL. of 0.1M NH3
solution. (of NH4OH = 4.8).
Given buffer is basic.
salt
pOH = + log
base
20 x0.1
22
= 4.8 + log 2 x 0.1
22
2
= 4.8 + log
0.2
= 4.8 + log 10 = 5.8
pH = 14 – pOH = 14 – 5.8 = 8.2.
Henderson’s equations for calculating the pH of an acidic buffer and pH of a basic buffer:
[ salt ]
For acidic buffer pH = P Ka + log
[acid ]
Where Ka is the ionisation constant of the weak acid
[salt]
pH = 14 - P K b - log For basic buffer
[base]
Where Kb is the ionisation constant of the weak base
Q. (a) Which is responsible for keeping the pH unchanged, even after adding a small amount of the
strong acid to the acidic buffer prepared by mixing acetic acid and sodium acetate solutions? Explain
why ?
(b) What is the pH of an acidic buffer prepared by mixing equal volumes of equimolar solutions of a
weak acid and its salt with a strong base?
(a) Acetate ion, because it removes the added H+ ions.
CH 3COO H CH 3COOH
[ salt ]
(b) pH = pKa , since pH = pKa + log
[acid ]
= pKa + log 1
pH = pKa + 0 = pKa
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Buffer capacity:
The no. of moles of a strong acid or a strong base required to change the pH value of one litre of buffer
solution by one unit is known as buffer capacity.
No.of moles of acid (or ) base added to onelitre of buffer solution
Buffer capacity =
[ P H change]
Applications of buffer solutions:
1) Buffer solutions are used to control PH of solution in chemical analysis, industrial synthetic reac-
tions and enzyme catalysed reactions.
2) They play a significant role in biochemical processes.
3) Buffer solutions are used in softening of hard water.
ACID – BASE INDICATORS:
Acid-base indicator:
“The external substance added to the contents of the titration flask in a visual detection method of end
point by a colour change” is known as an acid-base indicator.
Ex: Methyl orange, Phenolphthalein etc.
Use: To know the end point of the titration.
Q. Give the suitable indicators in different titrations and mention the pH rangeof different indicators.
Titration Indicator used pH change at the end point
1) strong acid + strong base Any Indicator like methyl orange, 3.3 - 10.7
Ex. (HCl + NaOH) methyl red, Phenolphthalein
2) strong acid + weak base methyl orange (or) methyl red 4.0 - 6.3
Ex. (HCl + NH4OH)
3) weak acid + strong base phenolphthalein 7.7 - 9.7
Ex. (CH3COOH + NaOH) .
4) weak acid + weak base no indicator is suitable as the pH The end point is not sharp.
Ex. (CH3COOH + NH4OH) change at the end point is not
sharp.
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pH ranges of common indicators :
Indicator pH range
Methyl orange 3.1 – 4.4
Methyl red 4.2 – 6.3
Litmus 5.0 – 8.1
Phenol red 6.8 – 8.4
Phenolphthalein 8.3 – 10.0
Q. (a) What is the suitable indicator in the titration of a CH3COOH and NaOH ?
(b) Which is the suitable indicator in the titration of a NH4OH and HCl ?
(c) Name the indicator which gives pink colour with a base and no colour with an acid?
(a) phenolphthalein (b) methyl Orange (c) phenolphthalein.
Q. Phenolphthalein indicator is not used in strong acid and weak base titration. Why?
The pH range of strong acid and weak base titration is 6.3 to 4.0 but the pH range of the phenol-
phthalein is 8.3 to 10.0. So phenolphthalein is not suitable for strong acid and weak base titration.
23. Methyl orange indicator is not used in weak acid Vs strong base titrations. Why?
The pH range of weak acid and strong base titration is 7.7 to 9.7. But the pH range of the methyl
orange is 3.1 to 4.4. So methyl orange is not suitable for weak acid and strong base titration.
Universal indicator:
An indicator that covers a very wide range of pH (3-11) and gives different colour changes at different
pH values is called as an “Universal Indicator”. It is used in determining the pH of a solution
Ex.: A mixture of phenolphthalien (0.1 gms), Methyl Red (0.2 gms), Methyl Yellow (0.3 gms),
Bromothymol Blue (0.4 gms), Thymol Blue (0.5 gms) dissolved in absolute alcohol and to which calculated
quantity of NaOH is added till an yellow coloured solution is formed. This gives different colours at different pH
values. They are
At pH = 2 Red At pH = 8 Green
At pH = 4 Yellow (Orange) At pH = 10 Blue
At pH = 6 Yellow
pH range of an Indicator:
“The range at which the colour of the indicator changes is called pH range of the Indicator”
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It depends on the value of the Indicator. Generally gives the Range.
It is used to know the pH change near the end points and choice of the Indicators used in different
volumetric titrations.
SALT HYDROLYSIS:
Salt hydrolysis:
The anion or cation or both of a salt react with water producing ions or H+ ions or both in
aqueous solution is known as salt hydrolysis.
salt + water base + acid
It is the reverse process of neutralisation.
Hydrolysis of sodium acetate:
Sodium acetate is the salt of strong base, NaOH and weak acid, CH3COOH. It ionises completely into
and Na+ ions.
CH3COONa CH 3 COO + Na+
CH3COO- is the strong conjugate base of the weak acid, CH3COOH. Hence is reacts with water
producing OH ions.
CH 3COO + H2O CH3COOH +
Therefore, the aqueous solution of sodium acetate is alkaline (basic). This is called anionic hydrolysis.
Kw
Hydrolysis constant ( K h ) = K
a
Hydrolysis of Ammonium Chloride (NH4Cl):
Ammonium Chloride is the salt of strong acid, HCl and weak base, NH4OH. It ionises completely to
form NH4+ and Cl- ions.
NH4Cl NH4+ + Cl-
NH4+ is the strong conjugate acid of weak base, NH3. Hence it reacts with water producing H+ ions.
NH4+ + H2O NH4OH + H+
Therefore, the aqueous solution of NH4Cl is acidic. This is cationic hydrolysis.
Kw
Hydrolysis constant ( K h ) = . K
b
Hydrolysis of Ammonium acetate (CH3COONH4):
Ammonium acetate is the salt of weak acid, CH3COOH and weak base NH4OH.
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It ionises completely into CH 3COO and NH4+ ions.
CH3COONH4 CH 3COO + NH4+
CH 3COO ion is the strong conjugate base of weak acid, CH3COOH and NH4+ ion is the strong conju-
gate acid of the weak base, NH3. Hence CH3COO- and NH4+ ions interact with water producing and H+ ions
respectively.
CH 3COO + H2O CH3COOH + OH-
NH4+ + H2O NH4OH + H+
H+ + OH- H2O
Kw
Hydrolysis constant ( K h ) = K K .
a b
As the anion and cation of the salt undergo hydrolysis , the aqueous solution of ammonium acetate is
neutral.
Expression indicating the relationship between the hydrolysis constant of a salt and the ionisation constants
of the weak acids and weak bases:
Type of salt hydrolysed Value of Type of Nature of the
hydrolysis Hydrolysis solution
constant (Kh)
1. salt of strong acid and Kw cationic hydrolysis acidic
weak base(Ex : NH4Cl) Kh
Kb
2. salt of weak acid and Kw anionic hydrolysis basic (or)
strong base (Ex : K h K alkaline
a
CH3COONa)
3. salt of weak acid and Kw both cationic and neutral
weak base (Ex K h K K anionic
a b
:CH3COONH4)
4. salt of strong acid and – – does not hydrolyse neutral
strong base (Ex : NaCl )
Q. Why is the solution of sodium bicarbonate alkaline?
The bicarbonate ion reacts with H+ ions of water forming unionised H2CO3 molecules and leaving
ions in solution.
+ H2O
H2CO3 + (anionic hydrolysis)
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Q. Which of the following hydrolyses? Why ? What is the nature of the solution resulting in cationic
hydrolysis?
i) NH4Cl ii) Na2SO4
NH4Cl hydrolyses, since it is the salt of weak base NH4OH and strong acid HCl.
NH4+ + H2O NH4OH + H+ (cationic hydrolysis)
Acidic [cation removes OH- ion from H2O leaving H+ ion]
Q. What is the nature of aqueous solutions of each one of the following?
i) KCN ii) FeCl3 iii) CuSO4 iv) Na2S v) NaCl vi) Na2CO3
a) What is the nature of the solution resulting in anionic hydrolysis?
b) What is the nature of the resulting solution involving both cationic and anionic hydrolysis?
i) alkaline ii) acidic iii) acidic iv) alkaline v) neutral vi) alkaline.
a) alkaline [ anion removes H+ ions from H2O leaving more in solution].
b) neutral or nearly neutral.
P. Calculate the hydrolysis constant (Kh) of a salt of sodium hydroxide and weak acid. If the Ka of the
acid is 2´10–6.
K w 1 1014
Kh = 0.5´10–8.
K a 2 106
P. Calculate the hydrolysis constant of a salt of weak acid (Ka= 2´10–6) and a weak base
(Kb = 5´10–7).
Kw 1 1014
Kh = 10–2.
K a K b 2 106 5 107
1. Calculate pH of the following solutions whose [H+] is
a) 6.2 x 10-4 b) 3.71 x 10 8 c) 3.6 x 10-10 d) 2 x 10-3
Ans: a) 3.22 b) 7.43 c) 9.4 d) 2.7
2. Calculate H+ ions concentration , OH- ion concentration and pH for the following solutions.
a) 0.5 M KOH b) 0.01M H2SO4 c) 0.001M Ba(OH)2
Ans: a) 13.7 b) 1.7 c) 11.3
3. For the following solution, calculate H+ ion concentration whose pH values are
a) 5.34 b) 2.3 c) 8.9
Ans: a) 4.57 x 10-6 b) 5 x 10-3 c) 1.26 x 10-9
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4. Find the pH of 0.2M CH3COOH and 0.5M NH4OH assuming complete ionisation.
(Ans: 0.6990, 13.7)
5. The pH of a buffer solution prepared by mixing 50mL. of 0.2M CH3COOH and 25mL of
CH3COONa is 4.8. What is the concentration of CH3COONa? (pKa = 4.8) (Ans: 0.4 M)
6. Calculate the pH value of the buffer containing 0.025 moles of NH4OH and 0.05 moles of NH4Cl
per litre. The dissociation constant of NH4OH is 1.8 x 10-5.
(Ans: 8.95)
7. 150ml. of 0.5N HCl and 100ml of 0.2N HCl are mixed. Find the pH of the resulting solution.
(Ans.: 0.42)
8. Equal volumes of 0.5N NaOH and 0.3N KOH are mixed in a experiment. Find the pOH and pH of
the resulting solution.
(Ans:13.6)
9. 50ml of 0.2M HCl is added to 30ml of 0.1M KOH solution. Find the pH of the solution?
(Ans: 1.06)
10. 40ml of 0.2N HNO3 when reacted with 60ml of 0.3M NaOH, gave a mixed solution. What is the
pH of the solution?
(Ans: 13.0)
11. The pH of an aqueous solution is 6.58. What are the concentrations of H+ ion and OH– ion.
(Ans: 3.8´10–8 M)
12. How many grams of NaOH are present per lit, if the pH is to be 10 ?
(Ans: 4´10–3 g)
13. What is the pH of 0.1M NaOH ?
(Ans: 13)
14. Find the pH of 10–8M HCl solution.
(Ans: 6.995)
15. 0.1825% HCl solution would have a pH of...........
(Ans: 1.301)
16. The pH of HCl solution is 5.4. What is the hydrogen ion concentration?
(Ans: 3.98´10–6 M)
17. 50ml of 0.2N H2SO4 were added to 100ml of 0.2N HNO3. Then the solution is diluted to 300ml.
What is the pH of the solution?
(Ans: 1.0)
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18. At a certain temperature, the ionic product of H2O is 9.55´10–14 moles2/lit2. Then what is the pH of
the solution?
(Ans: 6.51)
19. The Ka for a 0.2M CH3 COOH is 2´10–5. What is its pH ? (Assume CH3COOH as a strong acid)
(Ans: 0.6990)
20. Find the pH of 0.5M NH4OH solution. Kb for NH4OH is 2´10–5. (If NH4OH is assumed as a
strong base)
(Ans: 13.70)
21. What is the pH of HCl solution containing 3.65g is 250ml?
(Ans: 0.4)
22. A litre of buffer solution contains 0.1 mole of acetic acid and 1 mole of sodium acetate. Find the pH
if pKa of CH3COOH = 4.8.
(Ans: 5.8)
23. 20ml of 0.2M CH3COOH and 20ml of 0.4M sodium acetate were mixed together to form a buffer.
What is the pH. (CH3COOH has pKa of 4.8)
(Ans: 4.8)
24. Find the hydrolysis constant of 0.1M sodium acetate. Ka of CH3COOH = 2´105. (Ans : 5´10–10)
25. What is the pH of a buffer formed by mixing 20ml of 0.1M NH4Cl and 2ml of 0.1M NH3 ? pKb for
NH4OH = 4.8.
(Ans : 8.2)
26. The hydrolysis constant for NH4Cl solution is 0.5´10–9. Then what is the dissociation constant of
the base?
(Ans : 2´10–5)
27. The dissociation constant of NH4OH and CH3COOH are 2´10–5 and 1.8´10–5 respectively. Find
the hydrolysis constant of ammonium acetate.
(Ans: 2.8´10–5)
28. 50ml of 1M CH3COOH solution, when added to 50ml. of 0.5M NaOH, gives a solution with pH
value ‘X’. Find the value of ‘X’. (pKa of CH3COOH = 4.8)
March 2005 ) (Ans : 4.8)
29. What is the pH of a solution formed by mixing 50ml. 1M HCl and 50ml 0.1M NaOH ?
(Ans: 1.35)
30. The 0.005M monobasic acid has a pH of 5. What is the degree of dissociation?
(Ans : 0.2%)
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ASSIGNMENT
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1. Bronsted acid is a substance that
1) Loses electron pair 2) Donates proton
3) Accepts proton 4) Accepts electron pair
2. The conjugate base of H3O+ is
1) H+ 2) H2O 3) OH- 4) H2O-
3. The conjugate base of HSO 4 is
1) H2SO4 2) SO 4 2 3) SO4 4) H2O
4. Which of the following will behave as buffer solution ?
1) NaNO3 + HCl 2) NaNO3 + NH4OH
3) CH3COOH + CH3COONa 4) HCl + NaCl
5. The pH of 0.005 molar H2SO4 solution is
1) 2 2) 5 3) 10 4) 1
6. A certain buffer solution contains equal conc. of X- and HX. The Ka for HX is 10-10. The pH of buffer is
1) 4 2) 7 3) 10 4) 14
7. pH of a solution is given by the expression
1) log[H] 2) log1/[H+] 3) 1/logH+ 4) log/(H+)
8. The pH value of 0.001 M KOH solution is
1) 11 2) 3 3) 1 4) 4
9. The [H+] ion concentration in a given solution is 6´10-4. Its pH will be
1) 6 2) 4 3) 3.22 4) 2
10. What is the [OH-] concentration of a solution of pH = 3 ?
1) 10-11 mols/Lit 2) 10-3 mols/lit 3) 10-7 mols/lit 4) 10-14 mols/lit
11. Which of the following represent the conjugate acid base pair for the equilibrium reaction, HCl + H2O
H3O+ + Cl-
1) HCl, H2O 2) HCl, H3O+ 3) H3O+, H2O 4) Cl-, H2O
12. Which of the following is a Lewis base ?
1) AlCl3 2) SF4 3) SO3 4) C2H4
13. Conjugate acid of OH - base is
1) H2 2) H+ 3) H2O 4) H3O+
14. Which of the following is a Lewis acid ?
1) NaF 2) CaF2 3) BF3 4) CaCl2
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CHEMISTRY
15. NH 4 ion in aqueous solution will behave as
1) a base 2) an acid 3) neutral 4) both acid and base
16. Which of the following is an acidic salt?
1) NaHSO4 2) Na2SO4 3) Na2SO3 4) Na2S2O3
17. The conjugate acid of NH 2 is
1) NH 4 2) NH2OH 3) NH3 4) N2H4
18. The concentration of a weak acid (HA) is 0.1M. It’s degree of dissociation is 0.1%. Its dissociation
constant is
1) 1 ´ 10-3 2) 1 ´ 10-7 3) 1 ´ 10-10 4) 1 ´ 10-4
19. Which of the following behaves as amphiprotic substance ?
1) NH4OH 2) HCl 3) N2O 4) H2O
20. Which of the following is an aprotic solvent ?
1) Acetic acid 2) Water
3) Liquid ammonia 4) Carbon tetra chloride
21. The indicator suitable for the titration of CH3COOH with NaOH is
1) methyl orange 2) methyl red 3) phenolphthalein 4) None
22. Which of the following salts does not undergo hydrolysis when dissolved in water?
1) NaCl 2) NH4Cl 3) CH3COONH4 4) Na2CO3
23. Ammonia is considered to be a Lewis base because of
1) Polarity in molecule 2) presence of lone pair electrons
3) high volatility 4) the peculiar shape of the molecule
24. The pH of a solution containing 0.4g of NaOH in 1 lit solution is
1) 12 2) 2 3) 11 4) 13
25. Which of the following can act as an acid and also as a base?
1) H2SO4 2) HSO 4 3) SO 4 2 4) H+
26. In the reaction of water with NH3 given by the equation H 2 O + NH 3 NH + 4 + OH - water behave
as
1) An acid 2) A base
3) Both an acid and a base 4) Neutral
27. In the reaction BF3 + F - BF - 4 , BF3 can be considered as
1) Lewis Base 2) Lewis Acid 3) Lewis Salt 4) Bronsted acid
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28. A solution of pH 12.5 is
1) Neutral 2) Weakly basic 3) Strongly basic 4) Strongly acidic
29. In the system represented by the equation HS +OH S 2 + H2O, which of the following sub-
stances are Bronsted acids?
1) HS & H2O 2) HS & S 2 3) HS & OH 4) OH & H2O
30. In [Cu(NH3)4]SO4 complex, the Lewis acid is
1) SO42- 2) NH3 3) Cu+ 4) Cu2+
31. The [OH-] of a solution having a pH value 10 is
1) 10-10 moles/litre 2) 10-4 moles/litre
3) 3´10-4 moles/litre 4) 6´10-10 moles/litre
32. Which of the following salts undergo hydrolysis when dissolved in water?
1) NaCl 2) Na2SO4 3) NaNO3 4) Na2CO3
33. An aqueous solution whose pH is zero, is
1) neutral 2) acidic 3) alkaline 4) amphoteric
34. Which of the following indicators works in the pH range 8 to 9.5?
1) litmus 2) methyl orange 3) methyl red 4) phenolphthalein
35. The dissociation constants of acids HA, HB, HC and HD are 2.5´ 10-3, 7.5 ´10-5, 5.3 ´10-9 and
1.2 ´ 10-2 respectively. The weakest acid is
1) HA 2) HB 3) HC 4) HD
36. Aqueous solution of FeCl3 is acidic due to
1) cationic hydrolysis 2) anionic hydrolysis
3) hydrolysis of both the ions 4) dissociation
37. Solution of NaHCO3 is alkaline because of
1) bicarbonate ion 2) anionic hydrolysis
3) cationic hydrolysis 4) formation of NaOH
38. The following are some statements about pH:
I. It is log of reciprocal of [H+]
II. Acidic solutions have pH > 7 at 250C
III. Neutral solutions have pH = 0 at 250C
The correct statement/s is/are:
1) I, III 2) I only 3) II only 4) II, III only
39. Assertion (A): pH of 0.1 N H2SO4 solution is 1.
Reason (R) : A strong acid is completely ionised in aqueous solution
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CHEMISTRY
The correct answer is
1) Both A and R are true and R is correct explanation of A
2) Both A and R are true and R is not the correct explanation of A
3) A is true but R is false 4) A is false but R is true
40. Match List - I with List - II
LIST – I LIST – II
A. CH3COONH4 i. Lewis acid
B. AlCl3 ii. Salt of a weak acid and a weak base
C. NaHCO3 iii. Lewis base
D. NH3 iv. acidic salt
v. Arrhenius acid
The correct match is:
A B C D A B C D
1) ii i iv iii 2) iv i ii iii
3) iv ii iii i 4) ii iv i iii
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KEY
1 2 16 1 31 2
2 2 17 3 32 4
3 2 18 2 33 2
4 3 19 4 34 4
5 1 20 4 35 3
6 3 21 3 36 1
7 2 22 1 37 2
8 1 23 2 38 2
9 3 24 1 39 1
10 1 25 2 40 1
11 3 26 1
12 4 27 2
13 3 28 3
14 3 29 1
15 2 30 4
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4 SALTS
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Introduction
The term ‘salt’ is generally used in our daily life for common salt (sodium chloride) which is an important
component of our diet. In the language of chemistry, however, the term ‘Salts’ covers all electrovalent com-
pounds having positive and negative radicals formed by the reactions of corresponding acids and bases.
Therefore, the study of salts is invariably associated with the study of acids and bases. All salts are ionic
compounds and are electrically neutral as total positive charge on positive radicals is neutralised by negative
charge on the negative radicals. In their solid states, the oppositely charged ions are held together by strong
electrostatic forces of attraction. In a fused state or in aqueous solutions, these forces of attraction are broken
and these salts are dissociated into the ions. This accounts for the electrical conductivity of their solutions.
Thus, all salts are electrolytes in their molten states as well as in aqueous solutions. ‘Salts’ are defined as
the compounds formed by the replacement of ionisable hydrogen atoms of an acid either partially or wholly by a
metal or ammonium ion.
Example:
NaOH + H2SO4 ? NaHSO4 + H2O (Partial replacement)
2NaOH + H2SO4 ? Na2SO4 + 2H2O (Complete replacement)
Ionic definition of salt: Salt is an ionic compound which on dissociation in water yields a positive ion
other than H+ ion and a negative ion other than OH– ion.
Example:
MgCl2 ? Mg+2 + 2Cl–
(NH4)2SO4 ? 2NH4+ + SO4–2
Classification of salts
Salts are classified on the basis of their chemical composition into four types.
Simple salts: The salts which contain only one type of positive ion and one type of negative ion. Simple
salts are further classified into three types based on their chemical nature.
Comparative study of normal, acid and basic salts
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IIT - FOUNDATION - SET - V
Normal Salts Acid salts Basic salts
The salts which do The salts which
The salts which
not contain any contain one or
contain one or
Definition replaceable more
more replaceable
hydrogen or ‘OH’ replaceable
hydrogens.
groups. ‘OH’ groups.
By incomplete
By complete By incomplete
neutralisation
Formation neutralisation of neutralisation of
of acids and
acids and bases. acids and bases.
bases.
Partial
replacement of
All ionisable Partial replacement ionisable
Compositio hydrogens are of ionisable hydroxyl
n replaced by a Hydrogen by a groups by non
metal or NH4+ ion. metal or NH4+ ion. metal ions or
negative
radicals
Neutral in nature. Basic in nature.
Acidic in nature.
Properties Do not react with Reacts with
Reacts with bases.
acids or bases. acids
NaCl, (NH4)2SO4, NaHSO4, Cu(OH)NO3,
Al2(SO4)3 NH4HSO4 Cu(OH)Cl
e.g. e.g. e.g.
NH4OH + NaOH + H2SO4 Cu(OH)2 +
Examples H2SO4 HNO3
NaHSO4 + H2O
(NH4)2SO4 + Cu(OH)NO3+H
H2O 2O
Double salts: The salts which contain more than one simple salt are called double salts.
Example:
Potash alum ? K2SO4 · Al2(SO4)3 · 24H2O
Mohr’s salt ? FeSO4 · (NH4)2SO4 · 6H2O
Dolomite ? CaCO3.MgCO3
These salts are formed by the crystallisation and union of two simple salts dissolved in water.
Double salts undergo complete dissociation in aqueous solutions. They give reactions of all ions present in
both the simple salts.
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Example:
Mohr’s salt gives the reactions of NH4+, Fe+2 and SO4–2 ions in solution.
Note: These salts are usually associated with water of crystallisation.
Mixed salts: The salts which contain more than one basic or acidic radicals.
Example:
Sodium potassium carbonate (NaKCO3) (Na+, K+ and CO3–2 radicals)
Bleaching powder (CaOCl2) (Ca+2, OCl–, Cl– radicals)
Disodium potassium phosphate Na2KPO4 (Na+, K+, PO4–3 radicals)
These are formed by multi-step reactions involving corresponding acids and bases.
Mixed salts undergo complete dissociation in aqueous solutions. They give the reactions all the ions pro-
duced in aqueous solutions.
Example:
Bleaching powder gives reactions for Ca+2, OCl– and Cl– ions in solution.
Complex salts: The salts which contain one complex ion and one or more simple ions are called complex
salts.
Example:
Ag(NH3)]OH
K4[Fe(CN)6]
[Cu(NH3)](OH)2
These salts are formed by mixing saturated solutions of simple salts. This process is followed by
crystallisation.
A complex salt dissociates into a complex ion and simple ions. Simple ions react individually, while com-
plex ions react in groups.
Example:
K4[Fe(CN)6] Ì! 4K++[Fe(CN)6]–4
It gives reactions of potassium ions and ferro cyanide ions. But, it cannot give reactions of ferrous and
cyanide ions separately.
Preparation of salts
I.Preparation of soluble salts
1. By direct combination (synthesis): Heating a metal and a non-metal
together.
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Process: By passing chlorine gas over hot metal.
Salts prepared by this method: Chlorides of Aluminium, Iron and Zinc.
Examples:
(i) 2Al + 3Cl2 ? 2AlCl3
(ii) 2Fe + 3Cl2 ? 2FeCl3
2. By the action of dilute acids on metals: Active metals like Sodium,
Potassium etc., are treated with acids like HCl, H2SO4, HNO3 etc, to give
their respective salts.
Salts prepared by this method: MgSO4, ZnCl2, FeSO4, AlCl3
Examples:
(i) Zn + 2HCl ? ZnCl2 + H2?
(ii) Mg + H2SO4 ? MgSO4 + H2?
3. By the action of dilute mineral acids on oxides: All metal oxides are
insoluble and basic in nature. When the oxides are treated with dilute acids,
displacement takes place. Since metal oxides are basic in nature, neutralization
takes place.
Salts prepared by this method: FeSO4, Pb(NO3)2
Examples:
(i) FeO + H2SO4 ? FeSO4 + H2O
(ii) PbO + 2HNO3 ? Pb(NO3)2 + 2H2O
4. Action of dilute mineral acids on hydroxides:
(a) Soluble hydroxides are obtained by the titration of an acid with soluble
hydroxide.
Salts prepared by this method: (NH4)2SO4, NH4Cl
NaCl, Na2SO4, NaNO3
HCl, K2SO4, KNO3
Example:
2NH4OH + H2SO4 ? (NH4)2SO4 + 2H2O
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CHEMISTRY
(b) Insoluble hydroxide salts are obtained by the neutralisation of an acid and
its respective insoluble hydroxide.
Salts prepared by this method: FeSO4, Pb(NO3)2
Example:
Fe(OH)2 + H2SO4 ? FeSO4 + 2H2O
5. Action of dilute acids on carbonates
(a) Soluble carbonates salts are obtained by the titration of a carbonate salt
solution with an acid.
Salts prepared by this method NaNO3, (NH4)2SO4
Example:
(NH4)2CO3 + H2SO4 ? (NH4)2SO4 + H2O + CO2
(b) Insoluble carbonates salts are obtained by the neutralisation of an acid
with insoluble carbonate.
Salts prepared by this method FeSO4, Pb(NO3)2
Example:
PbCO3 + 2HNO3 ? Pb(NO3)2 + 2H2O + CO2
Qualitative analysis of soluble salts
For the identification of cations in the salts, NaOH and NH4OH are the generally used reagents. These
form insoluble precipitates with characteristic colours.
Table showing different characteristic precipitations given by different metal ions
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IIT - FOUNDATION - SET - V
Forme Colou
Catio Solubili
Precipita d r of Precipitat Precipitat
n of ty of
te NaOH NH4O NaO e NH4OH e NH4OH
salt NaOH
H H
Insolubl
Insoluble
Fe(OH) Dirty Dirty e in
Fe+2 Fe(OH)2 in excess
2 green Green excess
NH4OH
NaOH
Reddi Insolubl
Insoluble
Fe(OH) sh Reddish e in
Fe+3 Fe(OH)3 in excess
3 Brow Brown excess
NH4OH
n NaOH
Insolubl
No
Milky e in
Ca+2 Ca(OH)2 reactio
white excess
n
NaOH
Soluble
in
Soluble in
excess
excess
Zn(OH NaOH
Zn+2 Zn(OH)2 White White NH4OH
)2 (colour-
(colourles
less
s solution)
solution
)
Soluble
Chalk Insoluble
Pb(OH) Chalky in
Pb+2 Pb(OH)2 y in excess
2 White excess
White NH4OH
NaOH
Soluble in
Insolubl excess
Cu(OH Pale e in NH4OH
Cu+2 Cu(OH)2 Pale blue
)2 blue excess (deep
NaOH blue
solution)
Soluble
White White Soluble in
Al(OH) in
Al+3 Al(OH)3 gelati- gelatinou excess
3 excess
nous s NH4OH
NaOH
II. Preparation of insoluble Salts
1. By direct combination (synthesis): Heating the powdered metal with
sulphur.
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Salts prepared by this method: Zinc Sulphide, Iron Sulphide, Copper
Sulphide.
Example:
Fe + S ? FeS
2. Precipitation from soluble salts: In this method, when two soluble salt solutions are mixed,
double decomposition takes place and the insoluble salt is precipitated. This is the most common method for the
preparation of insoluble salts.
Salts prepared by his method: Carbonates, Sulphides, Sulphites, Sulphates, Chlorides.
Example:
BaCl2 + H2SO4 ? BaSO4 + 2HCl
Laboratory preparations of some specific salts
1. Iron chloride (FeCl3)
Method of preparation: Direct combination/Synthesis
Principle: Heated Fe powder reacts directly with dry Cl2 gas to form anhydrous FeCl3 which sublimes
and is collected by condensation.
Reaction: 2Fe + 3Cl2 ? 2FeCl3
Procedure:
Step 1: Dry Cl2 gas is passed through a combustion tube containing Fe
powder.
Step 2: As the reaction between Cl2 and Fe is exothermic, iron wires turn red
hot.
Step 3: FeCl3 volatilizes and condenses as brown scales. FeCl3 is highly
deliquescent, it is kept dry with the help of CaCl2 (absorbs moisture).
Precautions:
1. The apparatus should be air tight.
2. No external heating is required after some time.
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anhydrous
CaCl2
Fe wire
dry Cl2
Bunsen burner
conc.
H2SO4
solid Fecl3
2. Zinc sulphate ZnSO4 · 7H2O (White Vitriol)
Method of preparation: Displacement
Principle: Zinc displaces hydrogen from H2SO4 and forms ZnSO4 and it is crystallised later to get ZnSO4
· 7H2O.
Reaction: Zn(s) + H2SO4(aq) ? ZnSO4 + H2?
ZnSO4 + 7H2O ? ZnSO 4 .7 H 2O
( White vitriol )
Procedure:
Step 1: Dil.H2SO4 (1 vol. of acid and 5 vol of H2O) is taken in a beaker and
warmed, on a bunsen burner. Zn granules are added little by little while
stirring constantly.
Step 2: Effervescence takes place due to evolution of H2 (CuSO4 is added to
catalyse the reaction).
Step 3: Zn is added till all the acid has reacted (effervescence stops) and the
remaining metal settles at the bottom.
Step 4: The unreacted zinc which is in excess and black particles of carbon are
filtered and the filtrate is collected in a china dish
Step 5: The solution is evaporated till the crystals start appearing and it is then
cooled and filtered, washed and dried between the folds of a filter paper.
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White needle shaped crystals of hydrated zinc sulphate are formed.
3. Iron (II) sulphate FeSO4.7H2O (Green Vitriol)
Method of preparation: Displacement
Principle: Hydrogen from H2SO4 is displaced by Iron and form FeSO4. Later it is crystallised to get
FeSO4 · 7H2O.
Reaction: Fe(s) + H2SO4 ? FeSO4 + H2?
FeSO4 + 7H2O ? FeSO4 .7H2O
(Green vitriol)
Procedure:
Step 1: Dilute H2SO4 (1 vol. of acid and 5 vol. of H2O) is taken in a beaker and
heated on a bunsen burner.
Step 2: Iron filings are added little by little to it with constant stirring till the
effervescence stops (effervescence occurs due to H2 gas evolution).
Step 3: The hot solution is filtered and the filtrate is allowed to evaporate at room temperature.
The pale green colour crystals of FeSO4.7H2O are obtained.
4. Copper Sulphate: CuSO4 . 5H2O (Blue Vitriol)
Method of preparation Displacement:
Principle: Copper displaces Hydrogen from H2SO4 to form CuSO4. Later it is crystallised to get
CuSO4 · 5H2O.
Reaction: Cu(OH)2 + H2SO4 ? CuSO4 + 2H2O
CuO + H2SO4 ? CuSO4 + H2O
CuCO3 + H2SO4 ? CuSO4 + H2O + CO2?
CuSO4 + 5H2O ? CuSO4.5H2O
(Blue vitriol)
Procedure:
Step 1: Dilute acid is taken in a beaker and heated on a bunsen burner.
Step 2: Black CuO or CuCO3 is added little by little while stirring, till the excess settles down.
Step 3: It is filtered and its filterate is collected in a china dish and evaporated by heating it upto
crystallisation point and then allowing it to cool.
Step 4: The bright blue crystals of CuSO4.5H2O are collected and dried.
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Note
Similar preparation method can be used for other salts like MgSO4.7H2O,
Ca(NO3)2, CaCl2 etc
ZnO +H2SO4 ? ZnSO4 + H2O
PbO + 2HNO3 ? Pb(NO3)2 + 2H2O
CaO + 2HNO3 ? Ca(NO3)2 + H2O
Zn(OH)2 + H2SO4 ? ZnSO4 + 2H2O
MgCO3 + H2SO4 ? MgSO4 + H2O + CO2?
CaCO3 + 2HCl ? CaCl2 + H2O + CO2?
5. Sodium sulphate crystals Na2SO4.10H2O (Glauber’s salt)
Method of preparation: Neutralisation
Principle: Neutralisation of caustic soda with dilute H2SO4.
Reaction: 2NaOH + H2SO4 ? Na2SO4 + 2H2O
Na2SO4 + 10H2O? Na2SO4.10H2O
Procedure:
Step 1: 25cm3 of NaOH is transferred to a conical flask. A drop of
phenophthalein is added; it becomes pink.
Step 2: The acid (H2SO4 dilute) is taken in the burette and is titrated against
the alkali till the pink colour becomes colourless. This indicates that
the acid is completely neutralised by the base. At this moment titration
is stopped.
Step 3: The solution in the conical flask is transferred to the evaporation dish
and heated till the crystallisation point is reached. The dish is allowed
to cool till the crystals of Na2SO4.10H2O settle down.
Step 4: The crystals are filtered. The crystals obtained are then washed with
cold distilled water and dried in the folds of a filter paper to get the
crystals of Na2SO4.10H2O.
Note
This method is employed only to prepare salts containing Na+, K+ and NH 4 ions
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6. Preparation of PbCl2 and CaCO3
(i) Lead chloride (PbCl2)
Method of preparation: Double decomposition
Principle: Mixing the solutions of two appropriate compounds, each
containing one of the ions of the desired salt (salt required) results in the
formation of an insoluble solid ppt which can be filtered, washed and
dried.
Reaction: Pb(NO3)2 + 2HCl ? PbCl2? + 2HNO3
Pb(NO3)2 + 2NaCl ? PbCl2 ? + 2NaNO3
Procedure:
Step 1: Dil. HCl or NaCl is added to lead nitrate solution. A white
precipitate of PbCl2 appears.
Step 2: The heavy precipitate is filtered, washed and dried. Pure needle
shaped crystals of PbCl2 are obtained.
Note
Chlorides of Pb, Ag, Hg and sulphates of Ba, Pb, Ca can be prepared by this method.
(ii) Calcium carbonate – CaCO3
Reaction: CaCl2 + NaNO3 ? CaCO3 ? + 2NaCl
Procedure:
Step 1: Na2CO3 is added to a hot solution of CaCl2 in a beaker until Na2CO3 remains in excess.
Step 2: A white precipitate is formed by the interchange of radicals. It is filtered, washed and dried to
obtain an amorphous powder of Calcium Carbonate.
Note
Carbonates of all metals (expect Na, K, NH4+) can be prepared by this method.
General characteristics of salts
1. Physical state: Salts are non-volatile crystalline solids.
2. Electrical conductivity: Salts being electrovalent (ionic) compounds conduct current in molten
state as well as in aqueous state due to its dissociation into ions.
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3. Solubility: The salts are ionic compounds. But all the salts are not soluble in water. Some salts
remain insoluble.
Table showing the solubility of salts
Solubility of salts in water
Solubility of compounds Exceptions
All NH4 , Na , K compounds are soluble
+ + + --
All nitrates and nitrites are soluble --
All chlorides are soluble in water. Hg 2Cl2, AgCl, PbCl2
PbCl2 Soluble in hot water
All sulphates are soluble CaSO4, PbSO4, BaSO4
All oxides and hydroxides are insoluble Oxides of Na+, K+
All carbonates are insoluble. K2CO3, Na2CO3, (NH4)2CO3
All sulphides are insoluble. K2S, Na2S, (NH4)2S
Salt hydrolysis
The phenomenon in which a simple salt on dissolution in water forms a parent acid and a parent alkali
resulting in acidic, basic or neutral solution is called salt hydrolysis. It is the most important property of salts.
Depending on the nature of the salts, it is classified into four types.
1. Hydrolysis of salts of strong acids and weak bases
Example:
(NH4)2SO4, CuSO4, NH4NO3, FeCl3 etc.,
Hydrolysis of CuSO4:
CuSO4 + H2O ? Cu(OH)2 + H2SO4
CuSO4 is a salt and it exists in ionic form Cu(OH)2 is a weak base and it exists in molecular
form. H2SO4 is a strong acid and it exists in ionic form.
The ionic equation can be written as
Cu+2 + SO4–2 + 2H2O ? Cu(OH)2 + 2H+ + SO4–2
SO4–2 ions are spectator ions.
The net reaction can be
Cu+2 + 2H2O ? Cu(OH)2 + 2H+
Conclusion
(a) Cation of the salt reacts with water and hence it is called cationic
hydrolysis.
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(b) It results in the formation of a weak base and free H+ ions thereby
increasing the conc. of H+ ions in the solution.
(c) The pH values of the aqueous solution of these salts are less than 7 (0 to
7). The solutions are acidic in nature.
2. Hydrolysis of salts of weak acids and strong bases
Example:
Na2CO3, CH3OOK, CH3COONa etc.,
CH3COONa + H2O ? CH3COOH + NaOH
CH3COONa (Sodium acetate) being a salt exists in ionic form. CH3COOH is a weak acid so it exists in
molecular form. NaOH is a strong base hence it exists in ionic form.
The ionic equation can be written as
CH3COO– + Na+ + H2O ? CH3COOH + Na+ + OH–
Na+ ions are spectator ions.
The net reaction can be
CH3COO? + H2O ? CH3COOH + OH?
Conclusion
(a) Anion of the salt reacts with water and hence it is called anionic
hydrolysis.
(b) Results in the formation of weak acid and free OH– ions thereby increasing the concentra-
tion of OH– ions in the solution.
(c) The pH values of the aqueous solutions of these salts are more than 7 (7 t o 14). The solutions
are basic in nature.
3. Hydrolysis of salts of strong acids and strong bases
Example:
NaCl, K2SO4, KNO3 etc.,
Hydrolysis of NaCl:
NaCl + H2O ? NaOH + HCl
NaCl is a salt and it exists in ionic form. Both NaOH and HCl are strong
base and a strong acid respectively. Hence they also exist in ionic form.
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The ionic equation can be written as
Na+ + Cl– + H2O ? Na+ + OH– + H+ +Cl–
Both Na+ and Cl– are spectator ions
The net reaction can be
H2O ? H+ + OH–
Conclusions
(a) Neither cations nor anions of the salt react with water.
(b) These salts do not undergo hydrolysis as there is no net reaction
(c) As it produces both H+ and OH– ions in equal concentrations, the pH value i s
equal to 7. The solutions are neutral in nature.
4. Hydrolysis of salts of weak acids and weak bases:
Example:
(NH4)2CO3 , CH3COONH4
(NH4)2CO3 + 2H2O ? 2NH4OH + H2CO3
Ammonium carbonate is a salt and it exists in ionic form. NH4OH is a weak base and it exists in molecu-
lar form. H2CO3 is a weak acid and exists in molecular form.
Ionic equation can be written as
2NH4+ + CO3–2 + 2H2O ? 2NH4OH + H2CO3
There are no spectator ions. The net equation remains the same.
Conclusion
(a) Both cations and anions of the salt react with water.
(b) It results in the formation of a weak acid and a weak base. No free H+ or
OH– ions are formed.
(c) The pH value of the solution will be equal to 7. The solutions are neutral in n a -
ture.
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Amphoteric metals and their compounds
There are some metals which show characteristic properties of non metals in addition to metallic proper-
ties. Such metals are called amphoteric metals.
The most important amphoteric metals are zinc and aluminium.
Zinc and aluminium react with strong acids like HCl or H2SO4 giving corresponding salts with liberation of
hydrogen gas. Since this is a characteristic reaction of metals, these reactions indicate the metallic nature of Zn
and Al.
Zn + H2SO4 (dilute) ? ZnSO4 + H2?
2Al + 3H2SO4 (dilute) ? Al2(SO4)3 + 3H2?
Zinc and Aluminium give the same kind of reaction with strong alkalies. This shows the non-metallic prop-
erty of Zn and Al.
Zn + 2NaOH ? Na2ZnO2 + H2?
2Al + 2NaOH + 2H2O ? 2NaAlO2 + 3H2
From these reactions it can be concluded that Zn and Al are amphoteric in nature.
The oxides and hydroxides of these metals also give the same kind of chemical reactions, hence they are
also considered as amphoteric oxides and amphoteric hydroxides.
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5 ELECTRO CHEMISTRY
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CHEMISTRY
Conductors and non–conductors:
Substance, which allows the electric current to pass through it is called conductor.
Ex : Cu, Ag, graphite etc.
Substance which do not allow electric current to pass through it is called a non–conductor.
Ex : dry wood, rubber etc.,
Conductors are classified as metallic or electronic conductors and electrolytic conductors.
Classification of conductors:
i) Metallic (or electronic) conductors are the substances which conduct electricity through the move-
ment of electrons (from high negative potential region to lower positive potential region) without any
chemical changes.
Ex: Metals like Al, Cu, Ag and graphite etc.,
ii) Electrolytic conductors are the substances that conduct electricity through the movement of ions
towards oppositely charged electrodes accompanied by chemical changes.
Ex : Fused NaCl, Aqueous solutions of the salts.
Electrolytes and non-electrolytes:
Electrolytic conductors also called electrolytes.
A substance, that is, in the molten state or in the dissolved state containing ions and functioning as an
electrically conducting medium is called an electrolyte.
Ex : HCl, NaOH, KCl, AgNO3 etc.
A substance which even in the molten state or in the dissolved state does not contain ions and gives
a non–electrical conducting liquid is called non–electrolyte.
Ex : glucose, sugar, urea, alcohol etc.,
Electrolysis:
a) Electrolysis is the process of decomposition of chemical compound in the molten state or in the
solution state into its constituent elements under the influence of an applied electromotive force is
called electrolysis.
Ex : 2HCl (aq) Current
H2 + Cl2
b) When fused NaCl is electrolysed, sodium is formed at cathode and chlorine is evolved at anode.
NaCl Na + Cl
At Cathode: + e- Na (Reduction)
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At anode : (Oxidation)
Cl + Cl Cl2
Differences between electronic conductors and electrolytic conductors:
Electronic Conductors (Metallic Electrolytic Conductors
Conductors)
1. Electrons flow from negative end to 1. Ions moves towards oppositely
positive end. charged electrodes.
2. No new product is formed 2. New product is formed
3. Mass transfer is not involved 3. Mass transfer is involved
4. Conductivity decreases with increase 4. Conductivity increases with increase
of temperature of temperature
Q. What are the products obtained at electrodes when aqueous solution of Cu SO4 is electrolysed
between platinum electrodes? Give the electrode reactions.
Copper deposited on the cathode and oxygen at anode
Cu SO4 Cu+2 + (ionisation)
H2O H+ + OH- (very slightly ionized)
At cathode : Cu+2 + 2e- Cu
At anode 4OH- 2H2O + O2 + 4e-
Q. What products are obtained in the electrolysis of fused KCl and aqueous solution of KCl between
platinum electrodes?
Molten (Fused) KCl liquid undergoes electrolysis between ‘pt’ electrodes giving metal potassium
at the cathode and chlorine gas at the anode.
Cathode reaction : 2 K 2e 2 K
Anode reaction : 2Cl Cl2 g 2e
Electrolysis of aqueous KCl solution gives Hydrogen gas at the cathode and chlorine gas at the
anode.
Anode reaction: 2Cl
Cl2 2e deelectronation or oxidation
Cathode reaction: 4 H 2O 4e
4OH 2 H 2 electronation or reduction
But in the later stages of electrolysis oxygen is also obtained at the anode.
O2 4 H 4e
Cl Anode: 2 H 2O
ions are more easily oxidized than H2O molecules.
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CHEMISTRY
ELECTROLYSIS OF DIFFERENT ELECTROLYTES
S. Electrolyte Cathode Anode Products & electrode reaction Products & electrode
No. at cathode reaction at anode
H2 O2
1. Aq. K2SO4 Pt Pt
2 H 2 O 4OH
4 H 2 O _
4e
2 H 2 O
O2 4 H
4e
2. Aq. KCl Pt Pt H2 Cl2
2 e
2 H 2 O
H 2 2OH 2Cl
2e
Cl 2
3. Fused NaCl Pt Pt NA Cl
Na
e
Na e 1
Cl Cl 2
2
4. Fused Pt Pt Na O2
NaOH
4 Na 4 e
4 Na 4OH
4e
O2 2 H 2 O
5. Aq. CuCl2 Pt Pt Cu Cl2
2 e
Cu
Cu 2Cl
2e
Cl 2
6. Aq. CuSO4 Pt Pt Cu O2
4 e
2Cu
2Cu 2 H 2 O
O 2 4 H
4e
Faraday’s Laws:
Faraday’s Laws of electrolysis:
First Law: The mass of a substance deposited, liberated or dissolved at an electrode during electrolysis
of an electrolyte is directly proportional to the quantity of electricity passed through the electrolyte.
Where m = mass of the substance in grams and
Q = the quantity of electricity in coulombs. mQ
Where c = strength of current in amperes
t = time of passage of current in seconds m ct
Where e = constant called electrochemical equivalent. m = ect
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Second Law: When the same quantity of electricity is passed through solutions of different electrolytic
cells connected in series, containing different electrolytic solutions, the masses of different substances liberated or
deposited or dissolved at the electrodes are directly proportional to the chemical equivalents of the substances.
1m 1 E
m E
2 2
Where
m1 and m2 are the masses of the substances liberated at the electrodes and
E1 and E2 are their respective equivalent weights.
Electrochemical equivalent of an element:
Ans: The Electrochemical equivalent of an element is the mass of the element in grams liberated or depos-
ited by the passage of 1 coulomb of electricity through the electrolyte.
OR
The Electrochemical equivalent of an element is the mass of the element in grams liberated or depos-
ited by the passage of 1 ampere of current in 1 second through the electrolyte.
Relation between the Chemical equivalent and electrochemical equivalent:
Chemical equivalent of an element is the mass of the element in grams liberated or deposited at the
electrode by the passage of one Faraday (96500 c) of electricity through the electrolyte.
1 Faraday or 96,500 coulombs of electricity liberates or deposits E grams of the element, where ‘E’
is equivalent weight.
E
1 Coulomb of electricity liberates or deposits = grams
96500
E
Electrochemical equivalent, e = g / Coulomb.
96500
Atomic weight or formula weight
Chemical equivalent of an atom or ion =
Valency of the atom or ch arg e of the ion
Faraday:
The quantity of current that decomposes one gram equivalent of an electrolyte and deposits or
liberates one gram eq. wt. of the product at the electrodes is called Faraday. (OR)
The charge carried by 1 mole of electrons is called Faraday. its value is equal to 96, 5000 coulombs
1 Faraday = 96,500 Coulombs
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Charge carried by an electron = 1.66 x 10-19 coulombs
Charge carried by one mole of electrons = N x 1.66 x 10-19
= 6.02 x 1023 x 1.66 x 10-19 = 96500 C.
Problems:
P. A current of 0.25 amp. is passed P. A current of 10amp. is passed through
through CuSO4 solution for 45 molten AlCl3 for 96.5 seconds.
minutes . Calculate the amount of Calculate the mass of Al deposited at
copper deposited on the cathode (at. the cathode
wt of Cu = 63.6) (At. Wt. of Al=27)
Mct
A. m= m = mass. of Al =?
Mct zf
A: m
zf M = At. Wt. of Al = 27
M = 63.6 ; z = 2, F = 96, 500 coulombs Z = valency = 3
; C = 0.25 amp. F = 96, 500 coulombs
t = 45 mts. = 45 x 60 sec. C = 10 amp. ; t = 96.5 sec.
Wt of copper deposited, 27 10 96.5
63.6 m 0.09 g .
m= x 0.25 x 45 x 60 3 96500
2 x96500
= 0.224 g.
P. Calculate the current in amperes required P. 9.65 amp. Current is passed through
to deposit electrolytically 10g. of molten AlCl3 for one minute forty
Ag in 2 hours. (At. wt. of Ag = 108) seconds during electrolysis. The mass of
from aq. AgNO3 solution. Al deposited is 0.09g. at the cathode.
Mct What is the valency of Al?
A. m = m = mass. of Ag = 10g.
zf Mct
A. m = m = mass. of Al =0.09g.
M = At. Wt. of Ag = 108 zf
Z = valency = 1 M = At. Wt. of Al = 27
F = 96, 500 coulombs Z = valency = ?
C = ?. ; t = 2hrs. = 2 x 60 x 60 F = 96,500 coulombs
sec. C = 9.65 amp. ; t = 60+ 40 = 100 sec.
10 1 96500 27 9.65 100
C= 1.24amp. Valency 3
108 2 60 60 0.09 96500
Galvanic or Voltaic cells:
Galvanic (or) voltaic cell:
Volt meter
A galvanic or voltaic or electrochemical cell is a device which makes use of a spontaneous redox reaction
for the generation of electrical energy.
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IIT - FOUNDATION - SET - V
Salt bridge
Zn plate
anode Cu Plate
(-ve) cathode
(+ve)
ZnSO4(aq) CuSO4(aq)
Anode half cell Cathode half cell
Zn zn+2 + 2e– Cu+2 + 2e– Cu
Construction and working:
Construction and working:
Voltaic cell consists of two half cells i) anode half cell and ii) cathode half cell. The two half cells are
connected by a salt bridge containing a saturated solution of KNO3 in agar-agar gel. The anode half
cell consists of zinc rod dipped in ZnSO4 solution and the cathode half cell consists of a copper plate
dipped in CuSO4 solution. When the two half cells are connected externally through a volt meter,
current flows out from the cell due to potential difference.
The reactions taking place are:
2+ -
At anode: Zn Zn + 2e (Oxidation)
cathode: Cu2+ + 2e- Cu (reduction)
Net cell reaction : Zn(s) + Cu2+(aq.) 2+
Zn (aq.) + Cu(s)
During this chemical reaction current is produced as a result of conversion of chemical energy into
electrical energy. Daniel cell is a reversible cell because the cell reactions can be reversed when it is
connected to a battery of a higher potential.
Q. Write the two half cell reactions that take place in the following galvanic cells?
i) Pt, H2 (atm.) | HCl (0.1 M) | Cl2(atm.), Ptii) Zn|Zn2+ (1M) || Cu2+ (1M) |Cu
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CHEMISTRY
i) At anode H2 H+ + e (oxidation half cell reaction)
At Cathode : Cl2 + e- (Reduction half cell reaction)
Adding the two half cell reactions,
H2 + Cl2 H+ + ( cell reaction)
ii) At anode : Zn0 Zn2+ + 2 (oxidation half cell reaction)
At cathode: Cu2+ + 2e- Cu0 (Reduction half cell reaction)
Cell reaction, Zn0 + Cu2+ Zn2+ + Cu0
Salt bridge & its use:
Salt bridge is a ‘U’ – shaped glass tube filled with a saturated solution of KCl or KNO3 or NH4 NO3
in agar – agar gel. It is connected between two electrolytes of two half cells as a bridge in galvanic
cell.
Salt bridge is used to prevent the accumulation of charged ions around electrodes to maintain the
electrical neutrality in the two half-cells.
It prevents the physical mixing of two electrolytes present in the half-cells.
Electrolyte used in salt bridge:
KCl (or) KNO3 (or) NH4NO3
Purpose of a salt bridge in a galvanic cell:
To prevent the accumulation of charged ions around the electrodes and thus to maintain the electrical
neutrality in the two half-cells.
Single electrode potential:
The potential developed when a metal rod immersed in an aqueous solution of its salt or a gaseous
nonmetal in contact with the aqueous solution of its anion is called its single electrode potential.
Eg: Zn/ ZnSO4(aq.) or H2/ HCl(aq.)
Oxidation and reduction potentials:
The tendency of an electrode to gain electrons is called its reduction potential.
The tendency of an electrode to lose electrons is called oxidation potential.
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IIT - FOUNDATION - SET - V
Standard electrode potential:
The potential developed by an electrode (metal or non-metal) when it is in contact with its ions of
unit concentration at 250C is called its standard electrode potential. It is denoted by E0.
Ex : standard electrode potential of copper electrode is represented as
standard electrode potential of Cl2 electrode is represented as
Electrochemical series:
The arrangement of various electrodes in the increasing order of their standard reduction potentials
(SRP) is called an electrochemical series.
The electrode with highest negative SRP is placed at the top and the electrode with highest positive
value is placed at the bottom of this series.
Differences between electrolytic cell and galvanic cell:
Electrolytic cell Galvanic Cell
1. Electrical energy is converted to 1. Chemical energy is converted to electric
chemical energy. energy.
2. Anode is given +ve sign and cathode is 2. Anode is given -ve sign and cathode
given -ve sign. is given +ve sign.
3. Electricity is consumed 3. Electricity is produced.
4. Salt bridge or porous pot (diaphragm) is 4. Salt bridge or porous pot (diaphragm) is
not needed. needed.
Q. How do you construct a galvanic cell with the following electrodes?
Zn+2 (aq)/Zn, E0 = - 0.77V and Cd+2 (aq)/ Cd, E0 – 0.40 V
The electrode with lower reduction potential is taken as anode (L.H.E.)
The electrode with higher reduction potential is taken as cathode (R.H.E)
The cell is represented as,
Zn(S) | Zn+2 (aq) || Cd+2 (aq) | Cd(S)
EMF of a cell:
The difference between the standard electrode potentials of the electrodes in the two half cells is
known as EMF of the cell (Ecell).
The standard EMF of a Galvanic cell is calculated using the equation
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CHEMISTRY
E 0cell E 0cathode E 0anode (or)
0
E Cell = E 0RHE E 0LHE (or)
0
E Cell = E 0Re ducion electrode E 0Oxidation electrode
0
Ecell must be positive so that the reaction is spontaneous.
P.Calculate the of the following cells.
2
i) E0 of Zn2+(aq) | Zn = -0.77V and E0 of Cu aq | Cu = +0.38V
2
ii) E0 of Zn2+(aq) | Zn = -0.77V and E0 of Cd aq | Cd = -0.40V.
0 0 0
i) E 0cell = E cathode E anode ii) E cell = E 0cathode E 0anode
E 0cell = +0.38V – (-0.77V) = -0.40V – (-0.77V)
= -0.40V + 0.77V
= +0.38V + 0.77V
= +0.37 volts
= 1.11 volts
Q. Write the cell notation, oxidation half cell reaction, reduction half cell reaction, and the total cell
reaction for the cell constructed with the following electrodes.
Cu2+(aq) | Cu E0 = +0.34V
Ag+(aq) | Ag E0 = +0.80V
A. Cell notation is Cu|Cu2+(aq) || Ag+(aq) | Ag
Oxidation half cell reaction is Cu Cu2+ + 2e-
Reduction half cell reaction is 2Ag+ + 2e- 2Ag
Total cell reaction is Cu + 2Ag+ 2Ag + Cu2+
P. Construct the galvanic cell and calculate the EMF of the cell at 270C using the following two half cells.
0 0
A. ESn |Sn2 = +0.14V, ECd | Cd 2 = +0.40V
Taking reduction potentials
E0 0
ECd
Sn2 | Sn = -0.14V, 2
| Cd = -0.40V.
Construction of the cell : The electrode with more negative potential is to be made the anode and the
other cathode.
The cell is Cd | Cd 2 aq || Sn 2 aq | Sn
Anode Cathode
E 0cell E 0cathode E 0anode
= -0.14 – (-0.40)
= - 0.14 + 0.40 = +0.26volts.
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Nernst Equation:
Nernst equation &its significance:
Nernst equation relates the single electrode potentials with the concentration of the ions, temperature and
no. of electrons involved as follows:
2.303RT
0
i) Nernst equation for the cation or metal electrode is E E log M n
nF
Where R= 8.314Jk–1 mol–1, n= no. of electrons transferred. ; F= Faraday = 96500 C.
0.059
0
At 250C, E E log M n
n
2.303RT
0
ii) For anionic electrode, the nernst equation is E E log An
nF
0.059
0
Substituting the values for R, F and 250C as above E E log An
n
Significance: Nernst equation is useful to calculate the electrode potential of an electrode for any con-
centration of the ion from its standard electrode potential.
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CHEMISTRY
Problems:
P. Calculate the electrode potential of Zinc P. Calculate the electrode potential of the
electrode of Zn 2 = 0.1M single electrode.
0 A.
(E Zn 2 | Zn = -0.77v)
0.059
Cu aq C 0.01M | Cu ? E 0
0.337V
A: E = E0 +
n
. log M n
0.059
Zn2+ + 2e- Zn ( n=2) E E0 log C
n
0.059
E0 = -0.77V, n=2; M n = 10-1 E E0
2
log 10 2
0.059 E 0.337 0.03 2
E = -0.77 + log 10-1
2 E 0.337 0.06 0.277V
0.059 0.059
= -0.77 + x (-1) = -0.77 -
2 2
= -0.77 – 0.0295 E
Zn 2 ( 0.1M ) | Zn
= -0.7995V
P. Calculate the electrode potential of the P. Calculate the single electrode potential of
single
Pt , Cl 2 | Cl 0.01M ? E 0 1.36V
electrode Ag 0.01M | Ag ?
0.059
E 0 0.799V . E E0 log c
n
0.059
A. E E0 log c 0.059
n E E0 log 10 2
n
0.059
E E0 log 10 2 1
n 2 Cl 2 e Cl
0.059
0.799
1
log 10 2 Ag e Ag 0.59
1.36 2
= 0.799 0.059 2 1
0.799 0.06 2 = 1.36 0.12 1.48V
E = 0.799 – 0.120 = +0.679V
Q. Construct the galvanic cell from the following two half-cells and calculate the Ecell at 250C of the cell.
Fe | Fe2+ (aq) (1M) || Sn2+ (aq) (0.1M) | Sn
E 0 Fe2 |Fe = -0.44V
E 0 Sn2 |Sn = -0.14V
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IIT - FOUNDATION - SET - V
0.059
A: E Fe 2 |Fe = E0 + log [Fe2+]
n
Fe2+ + 2e- Fe ( n = 2)
Hence, n = 2, [Fe2+] = 1.0M; E0 = -0.44V
0.059
E Fe 2 (1M ) | Fe
= -0.44 + log [1M]
2
0.059
= -0.44V - x (0)
1
E Fe 2 |Fe = -0.44V
0.059
E = E0 + log [Sn2+]
Sn2 ( 0.1M ) | Sn n
Sn2++2e- Sn (n = 2)
n = 2; [Sn ] = 10-1M, E0 = -0.14V
2+
0.059
E = -0.14V + log [10-1M]
Sn 2 ( 0.1M ) | Sn 2
0.059
= -0.14V + x (-1)
2
= -0.14V – 0.0295V
= -0.1695V = -0.17V
E = -0.17V
Sn2 (0.1M ) | Sn
The cell therefore is, Fe | Fe2+ (1M)||Sn2+(0.1M) | Sn
Ecell = ESn - EFe
= - 0.17V – (-0.44V)
= - 0.17V + 0.44V
Ecell = +0.27V.
1
Q. Calculate the electrode potential of Zn Zn 0.1M Cl 0.01M Cl2 , Pt
2
E zn0 / zn 0.762V ; E1 =+01.36V
Cl / Cl
2
0.059
A. E zn / Zn = E 0
n
log Zn
0.059
E 0.762
2
log 10 1 0.762 0.03 1
E 0.762 0.03 0.792V
0.59 0.059
E1
Cl 2 / Cl
E0
1
log Cl 1.36
1
log 10 2
2
1.36 0.06 2 ; E 1.36 0.12 1.48V
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CHEMISTRY
1
P. Calculate the EMF of the cell Cu Cu 0.1M Cl 0.01M | Cl2 , Pt
2
0
E Cu
|Cu
0.337V ; E1 1.36V = +0.337–0.03=+0.307
Cl 2Cl
2
0.059
E1
E0 log 10 2 1.36 0.03 2 =1.36+0.06 2=1.36+0.12=1.48V
2
Cl 2 |Cl 1
0
E Cell E1 / 2 Cl E Cu =1.48–0.307 = +1.073V
2 |Cl |Cu
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ASSIGNMENT
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CHEMISTRY
1. Electrolysis of an aqueous solution of Sodium hydroxide between platinum electrodes yields.
1) Hydrogen at Cathode, Oxygen at Anode
2) Hydrogen at Anode, Oxygen at Cathode
3) Only Hydrogen at Cathode 4) Only Oxygen at Anode
2. Match the following
List – I List –II
A) Lithium 1) Charge on one mole of e-
B) Faraday 2) non–electrolyte
C) Urea 3) acts as anode only
D) KCl 4) Electrolytic conductor
5) Electrolyte in salt bridge
A B C D A B C D
1) 3 1 2 5 2) 3 1 5 4
3) 2 5 1 3 4) 3 5 4 2
3. When 1 ampere current is passed through a copper wire for 10 seconds, the number of electrons
passing through it is :
1) 1.6´1019 2) 1035 3) 1.6´1016 4) 6.2´1019
4. When 4 amperes of current was passed for 80 minutes through fused metal chloride, 4 grams of metal
was deposited at the cathode. The equivalent weight of the metal is (approx) :
1) 60 2) 40 3) 20 4) 32
5. Assertion (A) : Electrolysis of an aqueous solution of a bromide (or) Iodide gives Br2 (or) I2 at the
anode, where as that of NaF give O2 and not F2
Reason (R) : F2 is highly reactive where as Br2 and I2 are less reactive gases.
1) Both A and R are true. R is correct explanation of A
2) Both A and R are true but R is not the correct explanation of A
3) A is true R is false. 4) A is false R is true
6. During the electrolysis of Fused NaCl, the anodic reaction is
1) Oxidation of Cl- ions 2) Oxidation of Na+ ions
3) Reduction of Cl- ions 4) Reduction of Na+ ions
7. The passage of current through a solution of certain electrolyte results in the formation of Hydrogen at
cathode and Chlorine at anode. The electrolyte solution is:
1) Cupric Chloride in water 2) Sodium Chloride in water
3) Sulphuric acid in water 4) Water
8. A current of strength 2.5 amperes was passed through CuSO4 solution for 6 minutes 26 seconds. The
amount of Copper deposited at cathode is
1) 0.317 g. 2) 3.175 g. 3) 0.635 g. 4) 6.35g.
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IIT - FOUNDATION - SET - V
9. A certain quantity of electricity is passed through aqueous solutions of AgNO3 and Cupric sulphate
connected in series. The amount of Silver (at.wt.108) deposited at cathode is 1.08g. The amount of
Copper (at.wt.63.54) deposited is :
1) 0.6354 g. 2) 6.354 g. 3) 0.3177 g. 4) 3.177 g.
10. 36´103 coulombs of charge is passed through an electrolyte. If the current strength is 20 amperes, the
time for which current passed is :...............
1) 1 hr 2) 1/2 hr 3) 2 hr 4) 3/4 hr
11. When 96,500 coulombs of electricity is passed through a copper sulphate solution, the amount of
copper deposited at cathode will be :
1) 0.5 mole 2) 1.0 mole 3) 2.0 mole 4) 4.0 mole
12. Which of the following reactions occurs at anode during the electrolysis of copper sulphate solution
using copper electrodes :
1) Cu (s) Cu2+ (aq) + 2e- 2) Cu++ Cu (s) + 2e-
3) H+ + e- H 4) SO42- [aq] + 2H+ H2SO4
13. On electrolysis 1 mole Al atoms will be deposited by
1) 1 mole of elections 2) 2 moles of electrons
3) 3 moles of electrons 4) 6 moles of electrons.
14. A current of 2.5 Amperes passing through an electrolytic solution deposits 3 grams of metal in 50
minutes. The equivalent mass of the metal is :
1) 38.6 2) 51 3) 16.5 4) 74.6
15. What volume of chlorine at S.T.P. will be liberated at anode when one ampere current is passed for 32
minutes 10 seconds through HCl solution ?
1) 224 mL. 2) 448 mL. 3) 112 mL. 4) 896 mL.
16. The standard reduction potentials of x,y,z are respectively –0.126v, -3.045v and –0.25v. The strongest
reducing agent :
1) x 2) z 3) y 4) all are equally oxidants
17. Three moles of electrons are passed through three solutions in succession containing AgNO3, CuSO4
and AuCl3 respectively. The moles of cations reduced at cathode will be in the ratio of :
1) 1 : 2 : 3 2) 2 : 1 : 3 3) 3 : 2 : 1 4) 6 : 3 : 2
18. The volume of oxygen gas liberated at STP when 3.01´1024 electrons are flown through the circuit is :
1) 5.6 lit 2) 28 lit 3) 11.2 lit 4) 5.6 lit
19. Which of the following is a good conductor of electricity?
1) Aqueous solution of BaSO4 2) Aqueous solution of PCl3
3) Aqueous solution of HgCl2 4) HCl in Benzene
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CHEMISTRY
20. Mg(s) + Cd2+(aq) ® Mg2+(aq) + Cd(s), the cell is best represented by
1) Cd | Cd2+(aq)||Mg2+|Mg 2) Mg|Mg2+||Cd2+|Cd
3) Pt, Mg|Mg2+||Cd2+|Cd, Pt 4) Mg|Mg2+||Cd|Cd2+
21. The standard hydrogen electrode is represented as :
1) H 2|H+(1M) 2) Pt, H2(1atm.)|H+(1M)
3) H2(1 atm.) | H+, Pt 4) Pt, H+ | H2 (1 atm.)
22. For the cell Cu/Cu2+//Ag+/Ag, which of the following is true?
1) Cu electrode is anode 2) Cu electrode is +ve terminal
3) Ag electrode is anode 4) Ag electrode is -ve terminal
23. For the cell Zn/Zn2+//Cu2+/Cu, the cell reaction is :
1) Zn Zn2+ + 2 e -
2) Cu2+ + 2 e- Cu
3) Zn + Cu2+ Zn2+ + Cu 4) Zn2+ + Cu Zn + Cu2+
24. Which of the following statements is correct ?
1) Electrolysis of fused sodium hydride using pt electrodes yields H2 gas at Cathode
2) Electrolysis of aqueous K2SO4 using pt electrodes yields O2 gas at Cathode
3) Electrolysis of dilute HNO3 using pt electrodes yields O2 gas at Anode
4) Electrolysis of aqueous NaCl using pt electrodes yields H2 gas at Anode
25. Which one of the following dissolved in water forms a solution, which is non-conducting?
1) Epsom salt 2) Mohr salt 3) Urea 4) Chrome alum
26. A solution of Sodium Sulphate in water is electrolysed using inert electrodes. The respective products
at cathode and anode are :
1) H2, O2 2) O2, H2 3) O2,Na 4) O2, SO2
27. When standard zinc electrode is coupled with SHE :
1) Zinc electrode acts as anode 2) SHE acts as anode
3) Spontaneously reduction takes at zinc electrode 4) all the above
28. On passing a current through molten Aluminium chloride for some time, produced 11.2 L. of Cl2 at
N.T.P at anode, the quantity of aluminium deposited at cathode is:
1) 27 grams 2) 18 grams 3) 9 grams 4) 36 grams
29. The concentration of which one of the following electrolytes increases on electrolysis
1) HCl 2) NaCl(aq) 3)CuSO4(aq) 4) NaOH(aq)
30. Number of electrons passing per second through a cross section of copper wire carrying 10-6
amperes current is found to be : (approx)
1) 1.6´10-19 2) 6´10-35 3) 6´1012 4) 6´10-16
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IIT - FOUNDATION - SET - V
31. Which among the following yields hydrogen at anode on electrolysis ?
1) aqueous solution of NaCl 2) copper sulphate solution
3) fused NaCl 4) fused Sodium Hydride
32. The number of electrons passing per second through a cross section of silver wire carrying 965 am-
peres is nearly equal to
1) 6 10-21 2) 1 10+16 3) 6 10+29 4) 6.02 1021
33. The volume of Oxygen at N.T.P. liberated during electrolysis on passing 96500 coulombs of electricity
is :
1) 22.4 litres 2) 11.2 litres 3) 5.6 litres 4) 2.8 litres
34. The amount of Mg deposited by using one mole of electrons, during the electrolysis of fused MgCl2 will
be
1) 24 g. 2) 12 g. 3) 6 g. 4) 48 g.
35. The electrolyte used in salt bridge in a galvanic cell is :
1) Agar-agar 2) Gum Arabic 3) Gel 4) Potassium Nitrate
36. The SRP values of Ag/Ag+ and Zn/Zn2+ electrodes are 0.80 V and –0.76 V. In the cell built with these
two electrodes :
1) Ag electrode acts as anode and Zn electrode acts as cathode
2) Ag electrode acts as cathode and Zn electrode acts as anode
3) Both the electrodes act as cathode
4) the cell can’t be built with these two electrodes
37. When 9.65 coulombs of electricity is passed through a solution of silver nitrate, the amount of silver
deposited is :
1) 16.2 mg. 2) 21.2 mg. 3) 10.8 mg. 4) 6.4 mg.
38. One Faraday of electricity is passed separately through one litre of one molar aqueous solutions of : (i)
AgNO3 (ii) CuSO4 (iii) AuCl3 (iv) SnCl4. The ratio of number of moles of Ag, Cu, Au and Sn
deposited at cathode is
1) 1 : 2 : 3 : 4 2) 4 : 3 : 2 : 1 3) 4 : 2 : : 1 4) 3 : : 2 : 1
39. The electric charge required for deposition of 1 gram equivalent of substance at an electrode is :
1) 1 ampere/second 2) 9650 coulombs/seconds
3) Charge on 1 mole of electrons 4) 1 ampere/1 hour
40. The SRP value of Ni/Ni2+ electrode is –0.25 V. The emf of the cell built with Ni/Ni2+ and Fe/Fe2+
electrodes is 0.19 V. The SRP of Fe/Fe2+ electrode is :
1) 0.44 V 2) – 0.44V 3) 0.69V 4) – 1.9 V
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CHEMISTRY
KEY
1. 1 16. 3 31. 4
2. 1 17. 4 32. 4
3. 4 18. 2 33. 3
4. 3 19. 2 34. 2
5. 2 20. 2 35. 4
6. 1 21. 2 36. 2
7. 2 22. 1 37. 3
8. 1 23. 3 38. 3
9. 3 24. 3 39. 3
10. 2 25. 3 40. 2
11. 1 26. 1
12. 1 27. 1
13. 3 28. 3
14. 1 29. 4
15. 1 30. 3
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6 ANALYTICAL CHEMISTRY
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CHEMISTRY
ANALYTICAL CHEMISTRY
The process of identification of the chemical compounds with respect to their nature and quantity is called
‘Analysis’.
Analytical chemistry deals with the systematic identification of the various constituents of the chemical
compounds. Analysis of chemical compounds
Analysis of chemical compounds is basically classified into two types.
Chemical Analysis
Qualitative Quantitative
Analysis Analysis
Inorganic Organic Volumetric Gravimetric
Compounds Compounds Analysis Analysis
Quantitative analysis: The analysis of a compound with respect to the quantity of a substance or a constitu-
ent either in terms of weight or volume.
Qualitative analysis : The analysis of compound with respect to its chemical nature. The compound to be
identified may be a gas, a simple salt or an organic compound.
We deal with Qualitative analysis in this topic.
Qualitative analysis
(i) Identification of gases
Many compounds during chemical reactions evolve various gases. Based on the nature of the gas liberated
in a particular reaction, it is possible to predict the constituents of the respective compound. Hence, it is essential
to identify the nature of these gases based on their characteristic properties.
The common gases which we need to detect in chemical reactions are
(a) Oxygen (b) Hydrogen (c) Water vapour (d) Ammonia
(e) Chlorine (f) HCl (g) CO2 (h) SO2
(i) H2S (j) NO2
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IIT - FOUNDATION - SET - V
Name of the gas Characteristics of the gas
Oxygen (a) Colourless and odourless
(b) Does not respond to litmus test. It is neutral.
(c) Not inflammable as it can not catch fire on its own.
(d) Supporter of combustion. It keeps the burning splinter
glowing.
Hydrogen (a) Colourless and odourless gas.
(b) Neutral towards litmus.
(c) Inflammable as it catches fire.
(d) Burns with blue colour flame and givens out a pop sound.
(e) Not a supporter of combustion.
Water vapour (a) Colourless and odourless gas.
(b) Neutral to litmus.
(c) Not inflammable as it does not catch fire.
(d) Not a supporter of combustion.
(e) Condense to form tiny droplets of liquid when it comes in
contact with cooler parts of the container.
(f) Turns anhydrous CuSO4 from white to blue.
(g) Turns blue cobalt chloride paper to pink
Ammonia (a) Colourless gas with pungent odour.
(b) Produces tears on close contact.
(c) Turns moist red litmus blue. It is a basic gas.
(d) Fumes in moist air.
(e) Produces dense white fumes of Ammonium chloride
when a glass rod dipped in HCl is brought close to it.
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CHEMISTRY
Name of the gas Characteristics of the gas
Carbon dioxide (a) Colourless and odourless.
(b) Turns moist blue litmus red. It is an acidic gas.
(c) Not inflammable.
(d) Not a supporter of combustion.
(e) Puts off a glowing splinter.
(f) Turns lime water milky due to the formation of CaCO3.
Hydrogen chloride (a) Colourless gas with pungent odour.
(b) Turns moist blue litmus red. Acidic in nature.
(c) Forms white dense fumes of Ammonium chloride
when a glass rod dipped in NH4OH is exposed to it.
Hydrogen sulphide (a) Colourless gas.
(b) Has rotten egg smell.
(c) Turns moist blue litmus red. Acidic gas.
(d) Turns lead acetate paper black due to the formation
of lead sulphide.
Sulphur dioxide (a) Colourless gas with suffocating smell.
(b) Turns moist blue litmus red. Acidic gas.
(c) Bleaches the colour of litmus.
(d) Not inflammable.
(e) Not a supporter of combustion.
(f) Turns lime water white due to the formation of CaSO3.
(g) Turns potassium permanganate solution colourless.
It is a reducing agent.
(h) Turns orange coloured acidified potassium
dichromate solution to green.
Nitrogen dioxide (a) Reddish brown in colour.
(b) Pungent odour.
(c) Turns moist blue litmus red. It is an acidic gas.
(d) Turns moist Potassium Iodide brown. It is oxidizing
agent.
Chlorine (a) Greenish yellow in colour.
(b) Irritating smell.
(c) Turns moist blue litmus paper red and bleaches the
colour. Acidic gas and bleaching agent.
(d) Turns moist potassium bromide paper red. It is
oxidizing agent.
(e) Turns moist starch Iodide paper black.
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IIT - FOUNDATION - SET - V
(ii) Identification of simple salts
A simple salt contains an anionic part and a cationic part. The two radicals should be identified separately.
A. Identification of anions: The common anions which we come across in salts include carbonate, sulphide,
sulphite, chloride, nitrate and sulphate. These anions are basically divided into three categories on the basis of the
tests by which they can be identified.
(a) Dilute HCl group Carbonates, Sulphides and Sulphites.
These anions can be identified by their reaction with dilute acids.
(b) Concentrated H2SO4 group Chlorides and NitratesThese anions can be identified by their reaction
with conc. H2SO4.
(c) Barium chloride group Sulphates.
Sulphates can be identified in solution state by the addition of BaCl2 solution.
Anion Name of the Example of Observation
test Reactions involved
Colourless gas which
puts off burning
Action of splinter and turns lime
CaCO3 CaO + CO2
heat water milky as CO2 gas
is evolved.
Colourless gas with
CaCO3 + 2HCl brisk effervescence
Carbonate
Action of which puts off burning
dil. HCl H2O + CO2 splinter and turns lime
water milky as CO2 gas
is evolved.
(a) CaCO3 + BaCl2 CaCl2 A white precipitate is
Action of
+ BaCO3 formed.
BaCl2
(b) BaCO3 + HCl BaCl2 + The precipitate is
Solution
H2O + CO2 soluble in conc. HCl.
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CHEMISTRY
Anion Name of Examples of
Observation
the test` Reactions involved
Evolution of a
colourless gas with
(a) ZnS + 2HCl rotten egg smell.
ZnCl2 + H2S The gas gives a
Action of
Sulphide (b) (CH3COO)2Pb + H2S black precipitate
dil. HCl
when passed
2CH3COOH + PbS through lead
acetate solution as
the gas is H2S gas.
A colourless gas
with suffocating
Action of MgSO3 + 2HCl smell is evolved
dil. HCl MgCl2 + H2O + SO2 which turns
acidified K2Cr2 O7
green.
Sulphite
A white precipitate
(a) MgSO3 + BaCl2
Action of is formed. The
MgCl2 + BaSO3
BaCl2 precipitate is
(b) BaSO3 + 2HCl
Solution soluble in conc.
BaCl2 + H2O + SO2
HCl.
(a) BaCl2 + H2SO4 A colourless and
pungent smelling
BaSO4 + 2HCl gas is evolved. The
Action of
gas when exposed
conc.
(b) HCl + NH4OH to a glass rod
H2SO4
dipped in NH4OH
NH4Cl + H2O forms white dense
fumes. It is HCl gas.
Chloride
A greenish yellow
gas with strong
Action of ZnCl2 + MnO2 + H2SO4 pungent smell is
MnO2 and evolved.
conc. ZnSO4 + MnSO4 + Cl2+ The gas turns moist
H2SO4 H2 O potassium bromide
paper red.
It is chlorine gas.
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IIT - FOUNDATION - SET - V
Name of Examples of
Anion Observation
the test Reactions involved
A white curdy
BaCl2 + 2AgNO3 precipitate is
Action of
formed.
Chloride AgNO3
AgCl + Ba(NO3)2 The precipitate
Solution
dissolves in
NH4OH.
Action of Pb(NO3)2 + 2H2SO4
Reddish brown
hot conc.
fumes are evolved.
H2SO4 2PbSO4 + 2H2O + 2NO2
Action of Pb(NO3)2 + H2SO4 + Cu Reddish brown
Nitrate
Cu fumes are evolved
turnings more vigorously
and hot PbSO4 + Cu(NO3)2 + H2O and the solution in
conc. the test tube turns
H2SO4 blue.
A white crystalline
ZnSO4 + BaCl2 precipitate is
Action of
formed.
Sulphate BaCl2
ZnCl2 + BaSO4 The precipitate is
solution.
insoluble in conc.
HCl. It is BaSO4.
B. Identification of Cations
(a) Colour: Some of the cations can be identified by their colour, usually alkali and alkaline earth metal
cations are colourless. Some of the transition metal cations impart colours to their salts.
Cation Salt Colour
Cu+2 CuSO4.5H2O Blue
Cu(NO3)2 or CuCO3 Bluish
green
Fe+2 FeSO4 Light green
Fe+3 FeCl3, Fe2(SO4)3 Brown
MnCl2 Pink
Mn+2
MnSO4 Light pink
Except the above cations, the salts of all other cations are white in colour.
(b) Flame test
This is a specific test done for the detection of some cations which can impart colour of flame.
Procedure for the flame test
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CHEMISTRY
The metallic salt whose cation is to be detected is made into a paste with conc. HCl. This paste is exposed
to a non-luminous flame with a platinum wire. The colour of the flame changes depending on the cation present in
the salt.
Cation Colour of the flame
Sodium Golden yellow colour
Potassium Lilac colour
Zinc Green flashes
Barium Apple green
Calcium Brick red
Copper Light green
(c)Precipitation with sodium hydroxide \ ammonium hydroxide
Since almost all the metal salts give insoluble hydroxides with NaOH and NH4OH, these tests can be
used for identification of the cations in general. Some of the hydroxide precipitates are coloured which makes the
detection of the ions still easier. Some of the hydroxide precipitates are soluble in excess of alkali. By making use
of solubility of the precipitates in excess of reagents, the ions can be distinguished.
Addition of
Cation Observation (colour of precipitate)
reagent
Pb+2 + 2 NaOH Pb(OH)2 +2Na+
NaOH White
Addition of excess NaOH:
Pb(OH)2 + NaOH Na2PbO2 + 2H2O
Pb+2
Pb+2 + 2NH4 OH Pb(OH)2 + 2NH4+
NH4OH White
Addition of excess NH4OH:
No characteristic reaction.
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IIT - FOUNDATION - SET - V
Addition of
Cation Observation (colour of precipitate)
reagent
Cu+2 + 2 NaOH Cu(OH)2 + 2Na+
NaOH Light blue
Addition of excess NaOH:
No characteristic reaction
Cu+2
Cu+2 + 2NH4OH Cu(OH)2 + 2NH4+
NH4OH Light Blue
Soluble complex which is dark blue in colour.
Cu(OH)2 + 4NH4 OH [Cu(NH4)4] (OH)2 + 4H2O
Fe+2 + 2 NaOH Fe (OH)2 + 2Na
Light green
NaOH Addition of excess NaOH:
No characteristic reaction
Fe+2
Fe+2 + 2NH4OH Fe(OH)2 + 2 NH+4
NH4OH Light green
Addition of excess NH4OH:
No characteristic reaction
Fe+3 + 3NaOH Fe (OH)3 + 2Na+
Reddish brown
NaOH Addition of excess NaOH:
No characteristic reaction
Fe+3
Fe+3 + 3NH4OH Fe(OH)3 + 3 NH+4
Reddish brown
NH4OH Addition of excess NH4OH :
No characteristic reaction
Zn+2 + 2 NaOH Zn(OH)2 + 2 Na+
White
Addition of excess NaOH:
NaOH Zn(OH)2 + 2NaOH Na2ZnO2 + 4H2O
Soluble
Sodium
Zn+2 Zincate
Zn+2 + 2 NH4OH Zn(OH)2 + 2NH4+
White
Zn(OH)2 + 4 NH4OH [Zn(NH3)4] (OH)2 + 4H2O
NH4OH Soluble
Complex
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CHEMISTRY
Addition
Cation Observation (colour of precipitate)
of reagent
Al+3 + 3NaOH Al(OH)3 + 3Na+
White gelatinous
Addition of excess NaOH:
NaOH Soluble sodium aluminate is formed
Al(OH)3 + NaOH NaAlO2 + 2H2O
Al+3
Al+3 + 3NH4OH Al(OH)3+ 3NH4+
White gelatinous
NH4OH Addition of excess NH4OH:
Slightly soluble due to the formation of
[Al(NH3)(OH)4] complex.
Mn+2 + 2NaOH Mn(OH)2 + 2Na+
Flesh colour
NaOH Addition of excess NaOH
No characteristic reaction
Mn+2
Mn+2 + 2 NH4OH Mn(OH)2 + 2NH4+
Flesh colour
NH4OH
Addition of excess NH4OH:
No characteristic reaction
Mg+2 + 2 NaOH Mg(OH)2 + 2 Na+
Fine milky white
NaOH Addition of excess NaOH:
No characteristic reaction
Mg+2
Addition of NH4OH
NH4OH
No characteristic reaction.
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Detection of ammonium radical
Test Observation
1. Test with conc. NaOH: Colourless pungent smelling gas is evolved.
Concentrated NaOH It produces white dense fumes when a glass
solution is added to the rod dipped in HCl is introduced into the test
salt and the test tube is tube.
warmed. NH4Cl + NaOH NH3 + NaCl + H2O
NH3 + HCl NH4 Cl
2. Test with Nessler’s reagent: A brown precipitate is formed.
A few drops of Nessler’s
reagent is added to the salt
followed by the addition of
NaOH solution.
Action of heat on some simple salts
Some simple salts on dry heating give certain characteristic observations by which they can be identified.
Type of salt Observation
I Carbonate salts
(a) Except sodium, potassium
and ammonium carbonate, Evolution of carbondioxide gas.
all other carbonates
(b) Copper carbonate Light green amorphous powder changes to
black colour
CuCO3 CuO + CO2
(c) Lead carbonate White amorphous powder changes to reddish
brown.
Residue of (PbO) on cooling, changes to
yellow colour. Partly fuses in glass and stains
the test tube yellow.
PbCO3 PbO + CO2
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CHEMISTRY
Type of salt Observation
II Nitrate salts
(a) Lead Nitrate Evolution of reddish brown gas.
Cracking Sound
2Pb (NO3)2
2 PbO + 4 NO2 + O2
(b) Copper Nitrate Bluish green salt changes to black resisude of
CuO
2 Cu(NO3)2
2CuO + 4NO2 + O2
III Oxides
(a) Mercuric oxide Orange red amorphous powder changes to
deep red and finally to almost black
mercury.
2HgO 2Hg + O2
(b) Lead (II) Oxide Chocolate brown powder changes to reddish
brown, PbO. On cooling, residue changes to
yellow.
2 PbO2 2PbO +O2
(c) Trilead tetraoxide (Red Red powder changes to reddish brown.
lead) Residue on cooling turns yellow.
2Pb3O4 6PbO + O2
IV Hydrated salts- All hydrated salts
Lose water of crystallization.
Steamy vapours condense on the cooler
parts.
White crystalline salt swells and then
melts.
Na2CO3.10H2O Na2CO3+10H2O
Blue crystalline salt crumbles to form
white amorphous power.
CuSO4. 5H2O CuSO4+5H2O
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Type of salt Observation
Bluish green crystalline salt melts to give
bluish green mass.
Cu(NO3)2.6H2O Cu(NO3)2 + 6H2O
White crystalline salt changes to white
sticky mass.
Zn (NO3)2.6H2O Zn(NO3)2 + 6H2O
V Ammonium salts
(a) Ammonium Chloride The salt sublimes to form white dense
fumes.The fumes condense on cooler parts
of the tube forming white powdery mass.
NH4Cl NH3 + HCl
(b) Ammonium dichromate Orange red crystalline solid swells to a
large extent and gives off steaming fumes.
Residue is green.
(NH4)2Cr 2 O7 Cr2O3 + N2 + 4H2O
(Residue)
(c) Ammonium carbonate White amorphous salt with a strong smell
changes to opaque powder with the
evolution of Ammonia
(NH4)2 CO3 NH4HCO3 + NH3
VI Zinc salts
All zinc salts are white in colour. They
turn to yellow colour on heating and
again become white on cooling.
Zn(NO3)2 2ZnO + 4NO2 + O2
(Residue)
Nature of oxides and hydroxides of Metals
I. Metals and Metal oxides
All the metal oxides are basic in nature. They turn moist red litmus blue. They react with acids just like their
corresponding metals. The metals which are more electropositive than hydrogen liberate hydrogen gas on reacting
with acids.
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CHEMISTRY
Ex : Mg + H2SO4 MgSO4 + H2
MgO + H2SO4 MgSO4 + H2O
Some metals and their oxides are found to exhibit acidic nature also in addition to the usual basic nature.
Acidic nature is a characteristic property of non metals. Since these metals show both acidic nature and basic
nature they are called amphoteric metals and their oxides are called amphoteric oxides.
Ex : Zinc and Aluminium
Amphoteric Character of Zinc and Zinc Oxide
(a) Reaction with acid
Zn + H2SO4 ZnSO4 + H2
ZnO + H2SO4 ZnSO4 + H2O
Since Zn and ZnO react with acids in the above reactions, they exhibit basic nature.
(b) Reaction with bases
Zn + 2NaOH Na2ZnO2 + H2
ZnO + NaOH Na2ZnO2 +H2O
Amphoteric nature of Aluminium and Al2O3
Since Zn and ZnO react with base in the above reactions, they exhibit acidic nature.
(a) Reaction with acids
2Al + 3H2SO4 Al2 (SO4)3 + 3H2
Al2O3 + 3H2 SO4 Al2 (SO4)3 + 3H2O
The above reactions indicate the basic nature of Aluminium and Aluminium oxide.
(b) Reaction with bases
2Al 2 NaOH
2 H 2O
2 NaAlO 2 3H 2
Al2O3 + 2 NaOH 2 NaAlO2 + H2O
The above reactions of Aluminium and its oxide show their acidic nature.
Zinc and Aluminium can react with both acids and bases, so they are called amphoteric metals and their
oxides as amphoteric oxides.
Metal hydroxides
All metal hydroxides are basic in nature. They are called bases.
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IIT - FOUNDATION - SET - V
The bases which are soluble in water are called alkalis.
Ex : NaOH, KOH, NH4OH etc
Alkalis show all properties of bases. In addition, they also possess the property of solubility in water.
Among the common alkalis, calcium hydroxide shows least solubility in water.
Alkalis produce OH - ions in their aqueous solutions. These solutions are colourless. They can be identified
by litmus test and action of indicators.
Alkalis are classified into two types. They are a) strong alkalis b) weak alkalis.
(a) Strong alkalis
Strong alkalis ionise completely to give large concentration of OH- ions in their aqueous solutions.
Ex : NaOH, KOH
(b) Weak alkalis
Alkalis which ionize partially in their aqueous solutions and produce less concentration of OH- ions in the
solutions are called weak alkalis.
Ex: NH4OH, Ca(OH)2
Amphoteric hydroxides
Just like the metals and their oxides, the hydroxides of these metals are also amphoteric in nature. They
show both acidic and basic nature.
Reactions showing amphoteric nature
(a) Reactions with acids
Zn(OH)2 + H2SO4 ZnSO4 + 2H2O
2Al(OH)3 + 3H2SO4 Al2 (SO4)3 + 6H2O
(b) Reactions with bases
Zn(OH)2 + NaOH Na2 ZnO2 + 2H2O
Al(OH)3 + NaOH NaAlO2 + 2H2O
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336
MATHEMATICS
MATHEMATICS
Contents
1. RELATIONS REVISION AND
FUNCITIONS
2. REAL NUMBERS AND LIMITS
3. TRIGONOMETRY
4. MATRICES
5. FUNCTION
6. REMAINDER AND FACTOR
THEOREM
7. QUADRATIC EQUATIONS AND
INEQUATIONS
8. PROGRESSION
9. STATEMENTS
10. MATHEMATICAL INDUCTION
11 BINOMIAL THEOREM
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IIT - FOUNDATION - SET - V
RELATIONS REVISION
1 AND FUNCTIONS
338
MATHEMATICS
Function:
Let A and B be two non-empty sets. A relation f from A to B is said to be a function, if every element in A is
associated with exactly one element in B. It is denoted by f: A B (read as f is mapping from A to B). If (a,
b) f, then b is called the f image of a and is written as b = f(a). a is called the pre image of b.
Domain and Co-domain:
If f : A B is a function, then A is called domain and B is the co-domain of the function.
Range:
If f : A B is a function, then the set of all images of elements in its domain is called the range of f and is
denoted by f(A)
i.e., f(A) = {f(a) / aA}
Note: Range of a function is always subset of its co-domain i.e., f(A) B.
I. If f : AB is a function, and n(A) = m, n(B) = p, then the number of functions that can be defined from
A to B is pm
Example: A = {1, 2, 3, 4}; B = {2, 3, 4, 5, 6} are two sets. A relation f is defined as f(x) = x + 2. The relation f
: AB is a function and f = {(1, 3), (2, 4), (3, 5), (4, 6)}.
II. A = {–2, 2, 3, 4, }, B = {4, 9, 16} are two sets. The relation f, defined as f(x) = x2, is a function from A to B,
since every element in A is associated with exactly one element in B. The function is f = {(–2, 4), (2, 4),
(3, 9), (4, 16)}.
III. A = {–1, 1, 2, 5}, B = {1, 8, 125} are two sets. The relation f defined as
f(x) = x3 is not a function from A to B. The relation f = {(1, 1), (2, 8),
(5, 125)}. The number –1 is an element in A but it has no image in B.
IV. A = {1, 2, 3, 4}, B = {x, y, z, t, u, v} are two sets. A relation f is defined as follows.
f = {(1, x), (2, y) (3, z), (4, t) (1, u)}. Here f is not a function, because the element 1 in A is associated
with two elements x, u in B. Therefore, f is not a function.
Arrow Diagram:
Functions can be represented by arrow diagrams.
Example: A = {1, 2, 3, 4, 5}, B = {1, 4, 9, 16, 25, 36} are two sets.
A relation f is defined as function f(n) = n2
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IIT - FOUNDATION - SET - V
The arrow diagram of this function is given below.
Example: 1
A B
f
1 1
2 4
3 9
1
4
5 25
6
36
Every element in A is associated with exactly one element in B. So f : AB is a function.
Example: 2
A B
f
a
1
e
2
i
Every element in A is associated with exactly one element in B. So, f :A®B is a function.
Example: 3
A B
g
a
4
b
5
c
b is an element in A and it is not associated with any element in B. So g: A®B is not a function.
Example: 4
A B
h
2 4
3 9
16
4
2
5
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MATHEMATICS
2 is an element in A. It is associated with two different elements in B. i.e., 2 has two different images. So h : A
B is not a function.
Difference between Relations and Functions:
Every function is a relation but every relation need not be a function. A relation f from A to B is called a
function if
(i) Dom (f) = A,
(ii) no two different ordered pairs in f have the same first coordinate.
Example: Let A = {1, 2, 3, 4} B = {a, b, c, d, e}
Some relations f, g, h are defined as follows.
f = {(1, a), (2, b), (3, c), (4, d)}
g = {(1, a), (2, b), (3, c)}
h = {(1, a), (1, b), (2, c), (3,d), (4, e)}.
In the relation f the domain of f is A and all first coordinates are different. So f is a function. In the relation
g the domain of g is not A. So g is not a function. In the relation h the domain of h is A but two of the first
coordinates are equal i.e., 1 has two different images. So h is not a function.
Types of Functions:
One-one function or injection:
Let f : A B be a function. If different elements in A are assigned to different elements in B, then
the function f : A B is called a one-one function or an injection.
i,e., a1, a2 A and a1 a2 f(a1) f(a2) then f : A B is a one-one function (or) if f(a1) = f(a2)
a1 = a2 then f : A B is a one-one function.
Example: 1
A B
f
1 a
2 b
3 c
4 d
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Different elements in A are assigned to different elements in B. f : A B is a one-one function.
Example: 2
A B
g
1
5
2
6
4
1 and 2 are two different elements in A, but they are assigned to the same element in B. So g : AB is not
a one-one function.
Many to one function:
If the function f : AB is not one-one, then it is called a many to one function; i.e., two or more elements
in A are assigned to the same element in B.
Example: 1
A B
f
4
1
2
5
3
2, 3 are different elements in A, and they are assigned to the same element i.e., 5 in B. So f : AB is a
many to one function.
Onto function or surjection:
f : AB is said to be an onto function, if every element in B is the image of at least one
element in A.
i.e., for every b B, there exists at least one element a A, such that f(a) = b.
Note: If f : AB is an onto function, then the co-domain of f must be equal to the range of f. f(A)
=B
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MATHEMATICS
Example: 1
A B
f
1 a
2 b
3 c
4
Range = {a, b, c}
Co-domain = {a, b, c}
Range = co-domain
In the diagram, every element in B is the image of one element in A. Therefore. it is an onto function.
Example: 2
A B
1 a
2 b
3 c
d
In the above diagram, d is an element in B, but it is not the image of any element in A. Therefore, it is not
an onto function.
Into function:
If a function is not onto, then it is an into function, i.e., at least one element in B is not the image of any
element in A, or the range is a subset of the co-domain.
Bijective function:
If the function f : AB is both one-one and onto then it is called a bijective function.
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IIT - FOUNDATION - SET - V
Example: 1
A B
f
1 4
2 5
3 6
In the diagram, f : AB is both one-one and onto. So f is a bijective function.
Example: 2
A B
1 3
2 4
5
In the diagram, it is only one-one but not onto, so it is not bijective.
Inverse of a function:
If f : AB is function, then the set of ordered pairs obtained by interchanging the first and second coordi-
nates of each ordered pair in f is called the inverse of f and is denoted by f–1. i.e., f : AB is a function then its
inverse is f–1: BA
Examples:
(i) f = {(1,2), (2, 3), (3, 4)}
f–1= {(2,1), (3, 2), (4, 3)}
(ii) g = {(1,4), (2, 4), (3, 5), (4, 6)}
g–1 = {(4,1), (4, 2), (5, 3), (6, 4)}
(iii) A = {1, 2, 3, 4}, B = {x, y, z, a, b} are two sets the function h: AB is defined as follows.
h = {(1,a) (2, b) (3, x) (4, z)}
h–1= {(a,1) (b, 2) (x, 3) (z, 4)}
in the above examples only f–1 is a function, but g–1, h–1 are not functions
f : AB a function f–1 : BA need not be a function
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MATHEMATICS
Inverse function:
If f : AB is a bijective function then f–1: BA is also a function.
i.e the inverse of a function is also a function, only when the given function is bijective.
Example:
A = {1, 2, 3, 4. 5}; B = {a, e, i, p, u} A function f is defined as follows.
f = {(1, a), (2, e), (3, i), (4, u), (5, p)}. Clearly, f is a bijective function
Now f–1 = {(a, 1), (e, 2), (i, 3), (u, 4), (p, 5)}. Clearly f–1 is also a function and it is also bijective.
Identity function:
f : A A is said to be an identity function on A if f(a) = a for every a A it is denoted by IA
Example: A = {1, 2, 3, 4} The identity function on A is IA = {(1, 1), (2, 2),
(3, 3), (4, 4)}.
Note:
(i) Identity function is always bijective function.
(ii) The inverse of the identity function is the identify function itself.
Constant function:
A function f : A B is a constant function if there is an element b B, such that f(a) = b, for
all a A.
i.e., in a constant function the range has only one element.
Example:
A = {1, 2, 3, 4}; B = {a, e, i, u} are two sets and a function from A to B is defined as follows:
f = {(1, a), (2, a), (3, a), (4, a)}
f is a constant function.
Note: The range of a constant function is a singleton.
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Equal Function:
Two functions f and g, defined on the same domain D are said to be equal if
f(x) = g(x) for all x D.
Example:
1 Let f: R – {2} R be defined by f(x) = x + 2;
x2 4
and g: R – {2} R be defined by g(x) = ;
x2
f and g have the same domain R – {2}
Given f(x) = x + 2 ;
x 2 4 x 2 x 2
g (x) = = =x+2
x2 x2
f(x) = g(x) for all x R – {2}
Composite function:
If f : AB and g : BC are two functions, then the function g[f(x)] = gof from A to C, denoted
by gof is called the composite function of f and g.
In the composite function gof,
(i) the co-domain of f is the domain of g.
(ii) the domain of gof is the domain of f, the co-domain of gof is the co-domain of g.
(iii) composite function does not satisfy commutative property i.e gof fog.
(iv) if f : AB; g : BC ; h: CD are three functions, ho(gof) = (hog)of
i.e., the composite function satisfies associative property.
Real function:
If f : AB, and A and B are both subsets of the set of real numbers (R), then f is called a real
function.
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Graphs of functions:
A graph does not represent a function, if there exists a vertical line which meets the graph in two or more
points, i.e., a vertical line meets the graph at only one point, then the graph represents a function.
Example: 1
In this diagram, the vertical dotted line meets the graph at only one point. So the graph represents a
function.
Example: 2
In the above diagram, the vertical dotted line meets the graph at two points. So it is not a function.
Note:
(i) The y-axis does not represent a function.
(ii) The x-axis represents a many-one function.
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REAL NUMBERS
2 AND LIMITS
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MATHEMATICS
Real Numbers:
A number that can be represented on a number line is called a real number.
Rational Numbers (Q):
p
The numbers can be expressed in the form of (p, q Z, q 0) are called Rational numbers.
q
Irrational Numbers:
The numbers which are not rational numbers are called irrational numbers. i.e., The numbers that cannot
p
be expressed in the form where q is a non zero integer are called irrational numbers.
q
Example:
Integers: I or Z = { ……. –3, –2, –1, 0, 1, 2, 3,……}
Whole Numbers: W = {0, 1, 2, 3, 4, 5, 6, …….}
Natural Numbers: N = {1, 2, 3, 4, 5, 6,…….}
Properties:
(1) The square of a real number is always positive.
If m, n Z; a, b R (a 0)
1
(i) a–m = (ii) am.an = am + n
am
am n
(iii) n
am (iv)(am)n = amn
a
m
a am
(v) m
(ab) = a b m m
(vi) (b 0).
b bm
(vii) n a a 1/ n (viii) n a m a m / n
a na 1
(ix) n (x) m n a a mn
b na
(xi)a0 = 1 (a 0)
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Example:
Simplify (i) (27)5/3.
(27)5/3 = (33)5/3 = 35 = 243. ( (am)n = amn)
5 / 4
625
(ii)
16
5 / 4 5/ 4
625 16
= (a–m =
16 625
5/ 4
24 25 a m a m
= =
54 55 b bm
5
2
= .
5
MODULUS OF A REAL NUMBER (OR) Absolute value of a Number:
If x is a real number, then its absolute value or modulus value is denoted by x (read as mod x) and defined
as follows,
x = x when x > 0.
x = –x when x < 0.
x = 0 when x = 0.
Example:
The value of 7 7 ( 7 > 0)
The value of 5 (5) 5 ( –5 < 0)
The value of x is always positive.
Hence the Range of x is positive real numbers including ‘o’.
Properties:
(i) If = a then x = a.
(ii) If x a then –a x a.
(iii) If a then x –a or x a.
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MATHEMATICS
Example 1: If x = 8, then find the values of x.
Solution: We know that, If x = a; x = a
So x = 8, x = 8.
Example 2: If x 9, then find the range of x.
Solution: We know that, If x 9, then –a x a
x 9 –9 x 9.
Example 3: If x 15, then find the range of x.
Solution: If x a, the x –a or x a.
x 15 x –15 or x 15.
Example 4: If x 5 = 9, find the value of x.
Solution: We know that x = a, then x = a
x 5 = 9 x + 5 = 9.
i.e. x + 5 = 9 or x + 5 = –9
x = 4 or x = –14.
Example 5: 3x 7 = 22, find the value of x.
Solution: 3x 7 = 22 3x – 7 = 22
3x – 7 = 22 or 3x – 7 = –22
3x = 29 or 3x = –22 + 7
29
x= or 3x = –15 x = –5
3
29
x = or x = –5.
3
Example 6: Solve x 5 3 .
Solution: We know that x a –a < x < a
x 5 3 –3 < x–5 < 3.
–3 +5 < x < 3 + 5
2 < x < 8.
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Example 7: Solve 3x 10 43 .
Solution: We know that, x a x < –a or x > a.
3x 10 43 3x – 10 < –43 or 3x – 10 > 43.
3x < –43 + 10 or 3x > 43 + 10
53
3x < –33 or x > .
3
53
x < –11 or x > .
3
Graph of x :
Since x is any real number, x always positive, so graph of x belongs to first and second
quadrants only.
X O
X
Y
Limit of a Function:
Consider the following example:
f(x) = x2.
x 1.95 1.96 1.97 1.98
f(x) 3.8025 3.8416 3.8809 3.92
x 1.99 2.01 2.001 2.001
f(x) 3.96 4.04 4.004 4.00004
i.e., when x is close to 2, f(x) is close to 4.
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MATHEMATICS
Limit of a Function:
When x reaches to ‘a’, f(x) reaches to l, then we say the limit of f(x) when x tends to ‘a’ is l.
Lt f ( x ) = 1
It is denoted by x a
Properties:
(1) The limit of a function is unique.
(2) Lt [kf ( x )] k Lt f (x )
x a x a
(3) Lt [f (x ) g ( x )] Lt f ( x ) Lt g (x )
x a x a x a
(4) Lt [f (x ) . g( x )] Lt f ( x ) . Lt g (x )
x a x a x a
Lt f ( x )
f (x ) x a
(5) Lt Lt g( x )
x a g ( x )
x a
Example:
The limit of a chord of a circle is tangent to the circle.
If PQ is chord of a circle, when P approaches to Q along the circle, then the chord become a tangent
to the circle.
Problems:
While evaluating the limits we use the following methods.
(i) Method of Substitution.
(ii) Method of factorization.
(iii) Method of rationalization.
xn an
(iv) Using the formula xLt .
a x a
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(i) Method of Substitution:
In this method we directly substitute the value of x in the given function to obtain the limit value.
Examples: (1)
x 2
Lt x 2 5x 6 (2) 2 5(2) 6 20
x 2 4x 1 (3) 2 4(3) 1 22 11
(2) Lt
x 3
x 5
3 5 8 4
Indeterminate form:
Lt f ( x ) when x = a if f(a) = 0 form then we call it as an indeterminant form.
In x a 0
Method of factorization:
f (x)
If is indeterminate form, when x = a then there exist a common factor for f(x) and g(x). We remove
g( x )
the common factor, and use substitution method to find the limit.
x 2 5x 6
Example 1: Evaluate xLt
2 x 2 3x 2
4 10 6 0
Solution: When x = 2, 4 6 2 0 indeterminate form.
Now, x2 – 5x + 6 = (x–3)(x–2)
x2 – 3x + 2 = (x–1)(x–2)
x 2 5x 6 ( x 3) (x 2)
Lt 2 Lt
x 2 x 3x 2 x 2 ( x 1) ( x 2)
x 3 2 3
Lt 2 1 1
x 2 x 1
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MATHEMATICS
Method of rationalization:
f (x)
If is indeterminate form, and f(x), g(x) are irrational functions, then we rationalize f(x), g(x) and
g (x )
remove the common factor, then use substitution method to find the limit.
x
Example 1: Lt
Evaluate x 0 .
1 1 x
0 0
Solution: When x=0, 1 1 0 0 = indeterminate form.
Since g(x) is irrational form, so we multiply the numerator and denominator with rationalizing factor
of g(x).
x x (1 1 x )
Lt Lt
x 0 1 1 x x 0
x
= xLt (1 1 x ) = 1+
0 1 = 2.
(3 6 x )
Example 2: Evaluate Lt .
x 0 x3
3 63 0
Solution: When x=3, indeterminate form since f(x) is irrational form so we multiply the
33 0
numerator and denominator with rationalizing factor of f(x).
(3 6 x ) (3 6 x ) (3 6 x )
Lt Lt
x 3 x 3 x 3 x 3 (3 6 x )
9 (6 x )
= xLt
3 ( x 3)(3 6 x )
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3 x
= xLt
3 ( x 3)(3 6 x )
1 1 1 1
= xLt
3 (3 6 x ) (3 6 3 ) = (3 3)
= .
6
4 x 4x
Example 3: Evaluate xLt
0 9 x 9x
4 4 0
Solution: When x = 0, = it is indeterminate form.
9 9 0
Here f(x), g(x) both are irrational functions, we multiply both numerator and denomina
tor with their rationalizing factors.
4 x 4x 4 x 4x 9 x 9 x
Lt
x 0 9 x 9x 9 x 9 x 4 x 4x
(4 x ) (4 x ) 9 x 9x
Lt
x 0 (9 x ) (9 x ) 4 x 4x
2x 9 x 9x
Lt
x 0 2 x 4 x 4x
9 x 9x 9 9 2 9 3
Lt
x 0 4 x 4x 4 4 2 4 2
(iv) Using the formula:
xn an n 1
(i) x a x a na
Lt (where ‘n’ is a any rational number)
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MATHEMATICS
Proof:
Case 1: n is positive integer.
xn an ( x a ) ( x n 1 x n 2 a x n 3 a 2 ...... a n 1
Lt Lt
x a
x a xa xa
Lt x n 1 x n 2 a x n 3 a 2 .......... . a n 1 (n terms)
x a
= an – 1 + an – 2 a + an – 3 a2 + …….. + an – 1 (n terms)
= an – 1 + an – 1 + an – 1 + ……… + an – 1 (n terms)
= n.an – 1.
Case 2: n is negative integer.
Let n = –m (n is positive integer)
xn an x m a m
Lt Lt
x a x a x a x a
1 1 am xm
m
m m m
= Lt x a = Lt x a
x a x a x a xa
(x m a m )
= xLt
a (x a) x m a m
xm am 1
= xLt Lt m m
a x a x a x a
1
m a m 1 m ( m is positive integer using case(1))
a am
a m 1
m 2m
m a m 1 2 m
a
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xm an n 1
( m) a m 1 = xLt na (n = –m)
a x a
Case 3: n is rational number.
p
Let n where p, q are integers and q 0
q
p p 1 1
n
x a n x a q
q ( x ) (a ) p
q p q
Lt Lt Lt
= x a x a xa x a xa xa
Let x1/q = y then x = yq
a1/q = b then a = bq.
x a yq bq y b.
1 1
( x q ) p (a q ) p yp bp
Lt Lt q q
x a x a yb y b
yp bp
yb
Lt
y b yq b q
yb
yp bp
Lt
y b y b p .b p 1
yq bq q. b q 1 ( p, q are integers).
Lt
yb yb
p p
b p 1 ( q 1) b p q
q q
p p
p q( 1) p ( 1)
b q (b q ) q
q q
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MATHEMATICS
xn an n 1 p
Lt na ( n ; bq = a)
x a x a q
Hence when n is rational number,,
xn an n 1
Lt na
x a x a
xm a n m mn
Lt
Note: x a n n
a
x a n
x 4 256
Example 1: Lt ?
x 4 x 4
x 4 256 x 4 44
Solution: Lt Lt
x 4 x 4
x4 x 4
= 4 44 – 1
= 4 43 = 256.
x 5 243
Example 2:Evaluate xLt
3 x 2 9
x 5 243 x 5 35
Solution: Lt 2 Lt
x 3 x 9 x 3 x 2 3 2
5 5 2 xm an m mn
= 3 ( Lt n a )
2 x a x a n n
5 135
. 33 .
2 2
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x n 2n
Example 3: If xLt 32 find the value of n
2 x 2
x n 2n n 1
Solution: Lt n.2
x 2 x 2
32 = n.2n – 1
4 8 = n.2n – 1
4 23 = n.2n – 1
4 24 – 1 = n 2n – 1
n = 4.
Limits at x tend to infinity:
1 1
We know, when x tends to infinity tends to ‘0’. While evaluating limits at infinity put x = ;
x y
2x 3
Example 1: Evaluate xLt
x 5
.
1
Solution: Put x ; if x0, y0
y
2
3
2x 3 y
Lt Lt
x
x 5 y0 1 5
y
2 3y 2 0
Lt =
y0 1 5 y 1 0 = 2.
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MATHEMATICS
3x 2 4 x 5
Example 2: Evaluate xLt
4x 2 7
1
Solution: Put x = , when x , y0
y
1 1
3. 4. 5
3x 2 4x 5 y2 y
Lt Lt
x 4x 2 7 y0 1
4. 2 7
y
3 4 y 5y 2 3 0 0 3
Lt =
y0 4 7y 2 4 7(0) 4 .
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3 MATRICES
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MATHEMATICS
Introduction
We have learnt about the order of matrix, different kinds of matrices, some operations like transpose,
addition, subtraction, multiplication of matrices. We also learnt different properties of matrix multiplication.
In this chapter we will learn how to find determinant and inverse of a 2 x 2 square matrix. We also learn
about how to apply the concept of matrices to solve a system of linear equations in two variables.
Determinant:
a b
For a given 2 ´ 2 square matrix A = , the real number (ad – bc) is defined as the determinant of
c d
a b
A and is denoted by A or .
c d
2 5 2 5
Example: If A = , then determinant of A = A =
6 3 6 3
= 2(3) – (–5) 6 = 36.
Singular matrix:
If determinant of a square matrix is zero, then the matrix is called a singular matrix.
6 9
Example: For the square matrix A =
2 3
6 9
A = = 6 ´ 3 –9 ´ 2 = 18 – 18 = 0.
2 3
So A is a singular matrix.
Non–singular matrix:
If determinant of a square matrix is not equal to zero, then the matrix is called non–singular matrix.
2 4
Example: For the square matrix A = ,
5 3
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2 4
A = = 2(3) – (–4) 5 = 6 + 20 = 26 0.
5 3
So A is a non–singular matrix.
Multiplicative inverse of a square matrix:
For every non–singular square matrix A of order n, there exists a non–singular square matrix B of same
order, such that AB = BA = I. (Note that I is unit matrix of order n). Here B is called multiplicative inverse of A
and is denoted as A–1 B = A–1
1
Note: If AB = KI, then A–1 = B.
K
Multiplicative inverse of a 2 2 square matrix:
a b 1 d b
For a 2 2 square matrix A = , we can show that A–1 = c a
c d ad bc
1 d b
= A c a
Note:
1. For a singular square matrix A = 0, and so its multiplicative inverse doesn’t exist. Con-
versely if a matrix A doesn’t have multiplicative inverse, then A = 0.
1 –1
2. If A is a square matrix and K is any scalar, then (KA)–1 = A .
K
3. For any two square matrices A and B of same order (AB) –1 = B–1 A–1.
Method for finding inverse of a 2 2 square matrix.
a b
We know that for a square matrix A = .
c d
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MATHEMATICS
1 d b 1 d b
A–1 = = .
ad bc c a A c a
From this formula we can find A–1 using the following steps.
1. Find whether A =0 or not. If A = 0, then the given matrix is singular, so A–1 doesn’t exist.
If A 0, then the matrix has a multiplicative inverse and can be found by the following steps
(2), (3) and (4).
2. Interchange the elements of principal diagonal
3. Multiply the other two diagonal elements by –1.
1
Multiply each element of the matrix by
4.
A .
2 4
Example: Find the inverse of the matrix A =
3 5
2 4
Solution: A = –10 + 12 = 2 ¹ 0
3 5
A is non singular and A–1 exists.
1 d b 1 5 4
ad bc c a 2 3 2
A–1 = =
5 4 5
2 2 2 2
A–1 = 3 2 = 3
2 2 2 1
Solution of simultaneous linear equations in 2 variables.
The concept of matrices and determinants can be applied to solve a system of linear equations in two or
more variables. Here we present two such methods. First one is matrix inversion method and the second one is
Cramer’s method.
(1) Matrix inversion method:
Let us try to understand the method through an example.
Example: Solve the simultaneous linear equations
2x–5y= 1, 5x + 3y = 18.
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Solution:
Given system of linear equations can be written in matrix form as shown below.
2x 5y 1
5x 3y = 18
The L.H.S matrix can be further written as product of two matrices as shown below.
2 5 x 1
5 3 y = 18 or AX = B (1)
2 5 x 1
Here A = is called coefficient matrix, X = is called variable matrix and B = 18
5 3 y
is called constant matrix. Now we need to find values of x and y i.e., the matrix X
To find X premultiplying both sides of (1) with A–1.
A–1 (AX) = A–1 B.
or (A–1A) X = A–1 B [since A (BC) = (AB) C]
or I X = A–1 B, [A–1 A = I] or X = A–1 B, [IX = X]
X = A–1 B.
So to find X we have to find inverse of coefficient matrix (i.e. A) and multiply it with B.
1
x 2 5 1
X= = 18
y 5 3
1 3 5 1
= 5 2 18
2 3 (5) 5
1 3 5 1
=
6 25 5 2 18
1 3 1 5 18
= 5 1 2 18
31
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MATHEMATICS
1
93
31
1 3 90 1 93 3
= 31 = 31 = 1 =
31 31 31 1
31
x
Thus X = = x = 3, y = 1.
y
Thus in general any system of linear equations px + qy =a, and rx + sy = b can be
p q x a
represented in matrix form (i.e. AX=B) as = = b
r s y
p q x
Here A is coefficient matrix = , X is variables matrix = y and
r s
a
B = is constant matrix.
b
1
p q a
X = A–1 B =
r s b
Note:
1. Matrix inversion method is applicable only when the coefficient matrix A is non-singular i.e. A 0.
If A = 0, then A–1 doesn’t exist and so the method is not applicable.
2. This method can be extended for a system of linear equations in more than 2 variables.
(2) Cramer’s rule or Cramer’s method
This is another method of solving system of linear equations using concept of determinants. Unlike
matrix inversion method, in this method we don’t need to find the inverse of coefficient matrix.
Example: Solve the system of linear equations
3x + 4y = 2, 5x – 3y = 13 by Cramer’s method.
The system of equations can be written in matrix form (i.e. AX = B) as shown below.
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IIT - FOUNDATION - SET - V
3 4 x 2
5 3 y = 13
3 4 2
Here A = and B = 13
5 3
To find the solution by Cramer’s method we define two matrices B1 and B2. The matrix B1 is obtained by
replacing first column matrix A by the column in matrix B. similarly B2 is obtained by replacing column 2 of matrix
A by the column in B.
2 4 3 2
i.e. B1 = , B2 =
13 3 5 13
2 4
B1 13 3
Now x = = 3 4
A
5 3
2(3) 4(13)
=
3(3) 4(5)
6 52
= = 58/29 = 2 and
9 20
3 2
B2 5 13
y= = 3 4
A
5 3
3 13 2 5
= = 29/–29 = –1
3(3) 4 5
Thus in general for a system of linear equations px + qy = a,
rx + sy = b, solution by Cramer’s method is
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MATHEMATICS
a q p a
b s r b
x= p q ,y= p q
r s r s
Note:
1. If the coefficient matrix A is singular, then = 0, and so the method is not applicable.
2. This method can be extended to system of equations in more than 2 variables.
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4 TRIGONOMETRY
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MATHEMATICS
Earlier we have learnt about different trigonometric ratios, their relations and also the values ofthe trigono-
metric ratios for different acute angles. Here we shall learn how to find the signs of trigonometric ratios for all
kinds of angles.
Standard position of the angle:
The angle is said to be in its standard position if its initial side coincides with the positive X-axis.
Note:
(1) Rotation of terminal side in anti clock wise direction we consider the angle formed is positive and
rotation in clock wise direction the angle formed is negative.
terminal
X
O Initial
(2)Depending upon the position of terminal side we decide the angle in different quadrants.
Positive
O – X
negative
Coterminal angles:
The angles that differ by either 360° or the integral multiples of 360° are called coterminal angles.
Example:
60°, 360 + 60 = 420°, 2.360 + 60 = 780° are coterminal angles.
Note:
(1) If is an angle then its coterminal angle is in the form of (n .360 + )
(2) The terminal side of coterminal angles in their standard position coincides.
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360 + 60
60°
O O
Signs of trigonometric ratios:
(i) If lies in the first quadrant, i.e., 0 , then all the trigonometric ratios are taken positive.
2
(ii) If lies in the second quadrant, i.e., , then only sin and cosec are taken positive and all the
2
other trigonometric ratios are taken negative.
3
(iii) If lies in the third quadrant, i.e., ,then only tan and cot are taken positive and all the
2
other trigonometric ratios are taken negative.
3
(iv) If lies in the fourth quadrant, i.e., 2, then only cos and sec are taken positive and all the
2
other trigonometric ratios are taken negative.
Trigonometric ratios of (90° – ):
sin (90° – ) = cos ; cos (90° – ) = sin
tan (90° – ) = cot ; cot (90° – ) = tan
cosec (90° – ) = sec ; sec (90° – ) = cosec
Trigonometric ratios of (90° + ) :
sin (90° + ) = cos ; cos (90° + ) = – sin
tan (90° + ) = – cot ; cot (90° + ) = – tan
cosec (90° + ) = sec ; sec (90° + ) = – cosec
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MATHEMATICS
Trigonometric ratios of (180° – ):
sin (180° – ) = sin ; cos (180° – ) = – cos
tan (180° – ) = – tan ; cot (180° – ) = – cot
cosec (180° – ) = cosec ; sec (180° – ) = – sec
Trigonometric ratios of (180° + ):
sin (180° + ) = – sin; cos (180° + ) = – cos
tan(180° + ) = tan ; cot (180° + ) = cot
cosec (180° + ) = – cosec ; sec(180° + ) = – sec
Similarly, the trigonometric ratios of 270° and 360° can be written.
Note:
The trigonometric ratios of – are the same as the trigonometric ratios of 360° – .
So, sin (–) = sin (360° – ) = – sin and so on.
Examples:
1. Find the value of sin 225°.
Solution: 225° = sin (180° + 45°)
= – sin 45° = – 1
sin 225° = – 1
2. What is the value of tan 315°?
Solution: tan 315° = tan (360° – 45°)
= – tan 45° = – 1
tan 315° = – 1
3. Find the value of sin 150° + cos 210°.
Solution: sin 150° + cos 210° = sin (90° + 60°) + cos (180° + 30°)
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= cos 60° – cos 30°
1 3 1 3
= =
2 2 2
1 3
sin 150° + cos 210° =
2
4. Find the value of sin2 135° + sec2 135°.
Solution: sin 135° = sin(180° – 45°)
1
= sin 45° =
2
sec 135° = sec (180° – 45°)
= – sec 45° = – 2
2
1
sin 135° + sec 135° =
2 2 2
2
2
1 5
= 2
2 2
Heights and distances:
P1 B
P A
Let AB be a vertical line and PA and P1B be two horizontal lines as shown in the figure above.
Let APB = and PBP1 = . Then, (i) is called the angle of elevation of the point B as seen from
the point P and (ii) is called the angle of depression of the point P as seen from the point B.
Applications:
1. From a point on the ground which is at a distance of 50 m from the foot of the tower, the angle of
elevation to the top of the tower is observed to be 30°. Find the height of the tower.
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MATHEMATICS
Solution:
Let the height of the tower be ‘h’ metres.
AB 1 h
From PAB, tan 30° =
PA 3 50
50
h= or 50 3
3 3
50 3
Hence, the height of the tower is m
3
30°
P 50 m A
2. The angle of elevation of the top of a tower is 45°. On walking 20 m towards the tower along the line
joining the foot of the observer and foot of the tower, the angle of elevation is found to be 60°. Find
the height of the tower.
Solution:
Let the height of the tower be h metres.
Let QA = x metres.
In PAB,
AB h
tan 45° = 1=
PA 20 x
20 + x = h x = h – 20 – – (1)
From QAB,
AB
tan 60° =
QA
h
3 h= 3x h = 3 h 20 , (using (1))
x
3 1 h = 20 3
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h=
20 3
=
20 3 3 1
3 1
3 1 3 1
20 3 3
=
31
= 10 3 3
Hence, the height of the tower is 10 3 3 m.
B
h
meters
45° 60°
P 20 m Q A
3. From the top of a building 100 m high, the angles of depression of the bottom and the top of an
another building just opposite to it are observed to be 60° and 45° respectively. Find the height of the
building.
Solution:
Let the height of the building be h metres.
Let AC = BD = d metres
From BDE,
ED 100 h
tan 45° = 1=
BD d
d = 100 – h – – (1)
From ACE,
CE
tan 60° =
AC
100
3
d
3 d = 100
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MATHEMATICS
3 (100 – h) = 100 (using (1))
100
100 – h =
3
100
h = 100 –
3
3 1
= 100 =
100 3 3
3 3
Hence, the height of the tower is
100 3 3
m
3
E
45°
60°
45°
D B
d mts
h mts
60°
C d mts A
4. Two huts of equal height stand on either side of a road of 100m wide. At a point on the road and in
between the huts, the angles of elevation of their tops is found to be 60° and 45° respectively. Find
the height of each hut.
Solution:
Let the height of each hut be h metres
Let P be the point of observation on the road.
Let AP = x, Then PC = 100 – x
From PAB,
AB
tan 45° =
AP
h
1= x = h – – (1)
x
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IIT - FOUNDATION - SET - V
From PCD,
CD
tan 30° =
PC
h
3 3 (100 – x) = h
100 x
3 (100 – h) = h (using (1))
100 3 =
3 1 h
h=
100 3
=
100 3 3 1
= 50 (3 –
3 1 3)
3 1
Hence, the height of each hut is 50 (3 – 3 ) metres.
B D
h h
45° 60°
A x P 100 – x C
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MATHEMATICS
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5 FUNCTION
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MATHEMATICS
Function:
Let A and B be two non-empty sets. A relation f from A to B is said to be a function, if every element in A
is associated with exactly one element in B. It is denoted by f: A ® B (read as f is mapping from A to B). If (a, b)
Îf, then b is called the f image of a and is written as b = f(a). a is called the pre image of b.
Domain and Co-domain:
If f : A B is a function, then A is called domain and B is the co-domain of the function.
Range:
If f : A B is a function, then the set of all images of elements in its domain is called the range of f and is
denoted by f(A)
i.e., f(A) = {f(a) / aA}
Note: Range of a function is always subset of its co-domain i.e., f(A)B.
I. If f : AB is a function, and n(A) = m, n(B) = p, then the number of functions that can be defined
from A to B is pm
Example: A = {1, 2, 3, 4}; B = {2, 3, 4, 5, 6} are two sets. A relation f is defined as f(x) = x + 2.
The relation f : AB is a function and f = {(1, 3), (2, 4), (3, 5), (4, 6)}.
II. A = {–2, 2, 3, 4, }, B = {4, 9, 16} are two sets. The relation f, defined as f(x) = x2, is a function from
A to B, since every element in A is associated with exactly one element in B. The function is f = {(–
2, 4), (2, 4), (3, 9), (4, 16)}.
III. A = {–1, 1, 2, 5}, B = {1, 8, 125} are two sets. The relation f defined as
f(x) = x3 is not a function from A to B. The relation f = {(1, 1), (2, 8),
(5, 125)}. The number –1 is an element in A but it has no image in B.
IV. A = {1, 2, 3, 4}, B = {x, y, z, t, u, v} are two sets. A relation f is defined as follows.
f = {(1, x), (2, y) (3, z), (4, t) (1, u)}. Here f is not a function, because the element 1 in A is associated
with two elements x, u in B. Therefore, f is not a function.
Arrow Diagram:
Functions can be represented by arrow diagrams.
Example: A = {1, 2, 3, 4, 5}, B = {1, 4, 9, 16, 25, 36} are two sets. A relation f is defined as function f(n) = n2
The arrow diagram of this function is given below.
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Example: 1
A B
f
1 1
2 4
3 9
4 1
5 25
6
36
Every element in A is associated with exactly one element in B. So f : AB is a function
Example: 2
A B
f
a
1
e
2
i
Every element in A is associated with exactly one element in B. So, f :AB is a function.
Example: 3
A B
g
a 4
b
5
c
b is an element in A and it is not associated with any element in B. So g: AB is not a function.
Example: 4
A B
h
2 4
3 9
16
4
25
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MATHEMATICS
2 is an element in A. It is associated with two different elements in B. i.e., 2 has two different images. So h
: A B is not a function.
Difference between Relations and Functions:
Every function is a relation but every relation need not be a function. A relation f from A to B is called a
function if
(i) Dom (f) = A,
(ii) no two different ordered pairs in f have the same first coordinate.
Example: Let A = {1, 2, 3, 4} B = {a, b, c, d, e}
Some relations f, g, h are defined as follows.
f = {(1, a), (2, b), (3, c), (4, d)}
g = {(1, a), (2, b), (3, c)}
h = {(1, a), (1, b), (2, c), (3,d), (4, e)}.
In the relation f the domain of f is A and all first coordinates are different. So f is a function. In the relation
g the domain of g is not A. So g is not a function. In the relation h the domain of h is A but two of the first
coordinates are equal i.e., 1 has two different images. So h is not a function.
Types of Functions:
One-one function or injection:
Let f : A B be a function. If different elements in A are assigned to different elements in B, then the
function f : AB is called a one-one function or an injection.
i,e., a1, a2A and a1a2 f(a1) ¹ f(a2) then f : A B is a one-one function (or) if f(a1) = f(a2) a1 =
a2 then f : A B is a one-one function.
Example: 1
A B
f
1 a
2 b
3 c
4 d
Different elements in A are assigned to different elements in B. f : A B is a one-one function.
Example: 2
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IIT - FOUNDATION - SET - V
A B
g
1 5
2
6
4
1 and 2 are two different elements in A, but they are assigned to the same element in B. So g : AB is not
a one-one function.
Many to one function:
If the function f : AB is not one-one, then it is called a many to one function; i.e., two or more elements
in A are assigned to the same element in B.
Example: 1
A B
f
4
1
2
5
3
2, 3 are different elements in A, and they are assigned to the same element i.e., 5 in B. So f : AB is a
many to one function.
Onto function or surjection:
f : AB is said to be an onto function, if every element in B is the image of at least one element in A.
i.e., for every b B, there exists at least one element a A, such that f(a) = b.
Note: If f : AB is an onto function, then the co-domain of f must be equal to the range of f. f(A) = B
Example: 1
A B
f
1 a
2 b
3 c
4
Range = {a, b, c}
Co-domain = {a, b, c}
Range = co-domain
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MATHEMATICS
In the diagram, every element in B is the image of one element in A. Therefore. it is an onto function.
Example: 2
A B
1 a
2 b
3 c
d
In the above diagram, d is an element in B, but it is not the image of any element in A. Therefore, it is not
an onto function.
Into function:
If a function is not onto, then it is an into function, i.e., at least one element in B is not the image of any
element in A, or the range is a subset of the co-domain.
Bijective function:
If the function f : AB is both one-one and onto then it is called a bijective function.
Example: 1
A B
f
1 4
2 5
3 6
In the diagram, f : AB is both one-one and onto. So f is a bijective function.
Example: 2
A B
1 3
2 4
5
In the diagram, it is only one-one but not onto, so it is not bijective.
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IIT - FOUNDATION - SET - V
Inverse of a function:
If f : AB is function, then the set of ordered pairs obtained by interchanging the first and second coordi-
nates of each ordered pair in f is called the inverse of f and is denoted by f–1. i.e., f : AB is a function then its
inverse is f–1: BA
Examples:
(i) f = {(1,2), (2, 3), (3, 4)}
f–1= {(2,1), (3, 2), (4, 3)}
(ii) g = {(1,4), (2, 4), (3, 5), (4, 6)}
g–1 = {(4,1), (4, 2), (5, 3), (6, 4)}
(iii) A = {1, 2, 3, 4}, B = {x, y, z, a, b} are two sets the function h: AB is defined as follows.
h = {(1,a) (2, b) (3, x) (4, z)}
h–1= {(a,1) (b, 2) (x, 3) (z, 4)}
in the above examples only f–1 is a function, but g–1, h–1 are not functions
\f : AB a function f–1 : BA need not be a function
Inverse function:
If f : AB is a bijective function then f–1: BA is also a function.
i.e the inverse of a function is also a function, only when the given function is bijective.
Example:
A = {1, 2, 3, 4. 5}; B = {a, e, i, p, u} A function f is defined as follows.
f = {(1, a), (2, e), (3, i), (4, u), (5, p)}. Clearly, f is a bijective function
Now f–1 = {(a, 1), (e, 2), (i, 3), (u, 4), (p, 5)}. Clearly f–1 is also a function and it is also bijective.
Identity function:
f : A A is said to be an identity function on A if f(a) = a for every a A it is denoted by IA
Example: A = {1, 2, 3, 4} The identity function on A is IA = {(1, 1), (2, 2), (3, 3), (4, 4)}.
Note:
(i) Identity function is always bijective function.
(ii) The inverse of the identity function is the identify function itself.
Constant function:
A function f : A B is a constant function if there is an element bB, such that f(a) = b, for all a A.
i.e., in a constant function the range has only one element.
Example: A = {1, 2, 3, 4}; B = {a, e, i, u} are two sets and a function from A to B is defined as follows:
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MATHEMATICS
f = {(1, a), (2, a), (3, a), (4, a)}
f is a constant function.
Note: The range of a constant function is a singleton.
Equal Function:
Two functions f and g, defined on the same domain D are said to be equal if
f(x) = g(x) for all x D.
Example: 1
Let f: R – {2} R be defined by f(x) = x + 2;
x2 4
and g: R – {2} R be defined by g(x) = ;
x2
f and g have the same domain R – {2}
Given f(x) = x + 2 ;
x 2 4 x 2 x 2
g (x) = = x2 =x+2
x2
f(x) = g(x) for all x R – {2}
Composite function:
If f : AB and g : BC are two functions, then the function g[f(x)] = gof from A to C, denoted by gof is
called the composite function of f and g.
In the composite function gof,
(i) the co-domain of f is the domain of g.
(ii) the domain of gof is the domain of f, the co-domain of gof is the co-domain of g.
(iii) composite function does not satisfy commutative property i.e gof` fog.
(iv) if f : AB; g : BC ; h: CD are three functions, ho(gof) = (hog)of
i.e., the composite function satisfies associative property.
Real function:
If f : AB, and A and B are both subsets of the set of real numbers (R), then f is called a real function.
Graphs of functions:
A graph does not represent a function, if there exists a vertical line which meets the graph in two or more
points, i.e., a vertical line meets the graph at only one point, then the graph represents a function.
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IIT - FOUNDATION - SET - V
Example: 1
In this diagram, the vertical dotted line meets the graph at only one point. So the graph represents a
function.
Example: 2
In the above diagram, the vertical dotted line meets the graph at two points. So it is not a function.
Note:
(i) The y-axis does not represent a function.
(ii) The x-axis represents a many-one function.
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MATHEMATICS
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REMAINDER AND
6 FACTOR THEOREM
390
MATHEMATICS
Introduction
A real valued function f(x) of the form a0xn + a1xn–1 + . . . . . + an, (a0 0) is called as a polynomial
of degree n, where n is a non-negative integer. Here a0, a1, . . ., an are the coefficients of various
powers of x.
Examples:
(i) 4x6 + 5x5 + x4 + x2 – 1 is a polynomial in x of degree 6.
(ii) 2x3 + x2 + 1 is a polynomial in x of degree 3.
Note: A constant is considered to be a polynomial of zero degree.
In earlier classes we have learnt different operations on polynomials like addition, subtraction, multiplica-
tion and division. Here we shall learn two important theorems on polynomials.
Remainder theorem:
If p(x) is any polynomial and ‘a’ is any real number, then the remainder when p(x) is divided by (x – a) is
given by p(a).
Proof:
Let q(x) and r(x) be the quotient and remainder respectively, when p(x) is divided by x – a.
By division algorithm
Dividend = quotient divisor + remainder
i.e., p(x) = q(x) (x – a) + r(x)
If x = a, then
p(a) = q(a) (a – a) + r(a) r(a) = p(a)
i.e., p(x) = (x – a) q(x) + p(a)
Thus, the remainder is p(a)
Note:
1. If p (a) = 0. We say that ‘a’ is a zero of the polynomial p(x).
2. If p(x) is a polynomial and ‘a’ is a zero of p(x), then p(x) = (x – a) q(x).
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IIT - FOUNDATION - SET - V
b
3. If p(x) is divided by (ax + b), then the remainder is given by p
ba
4. If p(x) is divided by (ax – b), then the remainder is given by p
a
Example:
Find the remainder when the polynomial p(z) = z3 – 3z + 2 is divided by z – 2
Solution:
Given p(z) = z3 – 3z + 2
The remainder when p(z) is divided by z – 2 is given by p(2).
Now, p(2) = 8 – 6 + 2 = 4
= (2)3 – 3 (2) + 2
Hence, when p(z) is divided by (z – 2) the remainder is 4.
Factor theorem:
If p(x) is a polynomial of degree n ( 1) and a be any real number such that p(a) = 0, then (x
– a) is a factor of p(x).
Proof:
Let q(x) be the quotient and (x – k) (kR) be the factor of p(x)
Given p(a) = 0
By division algorithm
Dividend = quotient divisor + remainder
p(x) = q(x) (x – k) + p(a)
p(x) = q(x) (x – k) (5" p(a) = 0)
Therefore (x – k) is a factor of f(x), which is possible only if f(k) = 0
Hence (x – a) is a factor of p(x), (5"p(a) = 0)
Note:
1. If p(–a) = 0, then (x + a) is a factor of p(x).
b
2. If p = 0, then (ax + b) is a factor of p(x).
a
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MATHEMATICS
b
3. If p = 0, then (ax – b) is a factor of p(x).
a
4. If sum of all co-efficients of a polynomial is zero, then (x – 1) is one of its factor.
5. If sum of coefficients of oddpowers of x is equal to the sum of coefficients of even powers of x, then
one of the factors of the polynomial is (x + 1).
Examples:
(i) Determine whether (x – 3) is a factor of f(x) = x2 – 5x + 6
Given f(x) = x2 – 5x + 6
Now f(3) = (3)2 – 5 (3) + 6
= 9 – 15 + 6 = 0 f(3) = 0
Hence, by factor theorem we can say that (x –3) is a factor of f(x)
(ii) Determine whether (x – 1) is a factor of x3 – 6x2 + 11x – 6
Let f(x) = x3 – 6x2 + 11x – 6
Now f(1) = (1)3 – 6 (1)2 + 11 (1) – 6
= 1 – 6 + 11 – 6 = 0 f(1) = 0
Hence, by factor theorem we can say that (x – 1) is a factor of f(x)
Factorization of polynomials using factor theorem:
(i) Factorize x2 (y – z) + y2 (z – x) + z2 (x – y)
Let as assume the given expression as a polynomial in x say f (x)
f(x) = x2 (y – z) + y2 (z – x) +z2 (x – y)
Now put x = y in the given expression
f(y) = y2 (y – z) + y2 (z – y) +z2 (y – y)
= y3 –zy2 +y2z – y3 + 0 = 0 f(y) = 0
x – y is a factor of the given expression
Similarly if we consider the given expression as a polynomial in y we get
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IIT - FOUNDATION - SET - V
y – z is a factor of the given expression and we also get z – x is a factor of the
expression when we consider it as an expression in z.
Let x2 (y – z) + y2 (z – x) + z2 (x – y) = k(x – y) (y – z) (z – x)
For x = 0, y = 1, z = 2, we get
02 (1 – 2) + 12 (2 – 0) + 22 (0 – 1) = (0 – 1) (1 – 2) (2 – 0)
– 2 = – 2k
k=1
the factors of the given expression are x – y, y – z and z – x
(ii) Use factor theorem to factorize x3 + y3 + z3 – 3xyz
Given expression is x3 + y3 + z3 – 3xyz
Consider the expression as a polynomial in variable x say f(x)
i.e., f(x) = x3 + y3 + z3 – 3xyz
Now
f [–(y + z)] = [– (y + z)]3 + y3 + z3 – 3 [–(y + z)] yz
– (y + z)3 + y3 + z3 + 3yz (y + z)
= – (y + z)3 + (y + z)3 = 0 f [–(y + z)] = 0
according to factor theorem x – [– (y + z)] i.e., x + y + z is a factor of f(x)
i.e., x3 + y3 + z3 – 3xyz
Now we use the long division method to get the other factor as
x2 + y2 + z2 – xy – yz – zx
x3 + y3 + z3 – 3xyz = (x + y + z) ( x2 + y2 + z2 – xy – yz – zx)
Horner’s Process for synthetic division of Polynomials:
When a polynomial f(x) = p0xn + p1xn – 1 + - - - + pn – 1x + pn is divided by a binomial x – , let
the quotient be Q (x) and remainder be R.
We can find quotient Q(x) and remainder R by using Horner’s synthetic division process as
explained below.
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MATHEMATICS
po p1 p2 ---- pn–1 Pn --- Ist row
left (corner)
q0 q1 ----qn – 2 qn–1 --- IInd row
qo q1 q2 ---- qn–1 R --- IIIrd row
Step 1:
Write all the coefficients p0, p1, p2, - - - - - , pn of the given polynomial f(x) in the order of descending
powers of x as first row. When any term in f(x) (as seen with descending powers of x) is missing we write zero for
its coefficient.
Step 2:
Divide the polynomial f(x) by (x – ) by writing a in the left corner as shown above (x – = 0 x =)
Step 3:
Write the first term of the third row as q0 = p0 then multiply q0 by a to get q0 and write it under p1, the first
element of the second row
Step 4:
Add q0 to p1 to get q1, the second element of the third row
Step 5:
Again multiply q1 with to get q1 and write q1 under p2 and add q1 to p2 to get q2 which is the third
element of the third row
Step 6:
Continue this process till we obtain qn – 1 in the third row. Multiply qn – 1 with and write qn – 1 under pn and
add qn – 1 to pn to get R in third row as shown above
In the above process the elements of the third row i.e., q0, q1, q2, - - - , qn–1 are the coefficients of the
quotient Q(x) in the same order of descending powers starting with xn – 1
Q (x) = q0x n – 1 + q1x n – 2 + - - - - + qn – 2 x + q n – 1
and the remainder R i.e., the last element of the third row
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Note:
If the remainder R = 0 then a is one of the roots of f(x) = 0 or x – is a factor of f(x)
Example:
Factorize x4 – 10x2 + 9
Let p(x) = x4 –10x2 + 9
Here sum of coefficients = 0, and also
sum of coefficients even powers of x = sum of coefficients of odd powers of x
(x – 1) and (x + 1) are the factors of p(x).
Multiplier of x – 1 is 1 and x + 1 is –1
The quotient is x2 – 9
Hence p(x) = (x – 1) (x + 1) (x2 – 9)
p(x) = (x – 1) (x + 1) (x – 3) (x + 3)
Problems based on factor and remainder theorems:
1. Find the value of a if ax3 – (a + 1) x2 + 3x + –5a is divisible by (x –2).
Solution:
Let p(x) = ax3 – (a + 1) x2 + 3x + –5a
If p(x) is divisible by (x – 2), then its remainder is zero i.e., p(z) = 0
a(2)3 – (a + 1) (2)2 + 3(2) –5a = 0
8a – 4a – 4 + 6 – 5a = 0
–a+2=0
a=2
The required value of a is 2.
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MATHEMATICS
2. If the polynomial x3 + ax2 + bx –30 is exactly divisible by x2 – 2x – 15. Find a, b and also the third
factor.
Solution:
Let p(x) = x3 + ax2 + bx – 30
Given p(x) is exactly divisible by x2 – 2x – 15 i.e., (x – 5) (x + 3)
p(x) is divisible by (x + 3) and (x – 5)
p(–3) = 0 and p(5) = 0
Consider p(–3) = 0
(–3)2 + a(–3)2 – b(–3) –30 = 0
–27 + 9a + 3b – 30 = 0
9a + 3b – 57 = 0
3a + b – 19 = 0 (1)
Now consider p(5) = 0
i.e., 53 + a (5)2 – b(5) – 30 = 0
125 + 25a – 5b – 30 = 0
25a – 5b + 95 = 0
5a – b + 19 = 0 (2)
Adding (1) and (2), we get
3a + b – 19 = 0
5a – b + 19 = 0
8a = 0 a = 0
Substituting a in (1), we get
3(0) + b – 19 = 0 b – 19 = 0 b = 19
The required values of a and b are 0 and 19 respectively
p(x) = x3 + 0(x2) – 19x –30
i.e., p(x) = x3 – 19x –30
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Thus, the third factor is x + 2.
3. Find a quadratic polynomial which when divided by (x – 3) and (x + 1) leaves remainders 2 and 3
respectively and is exactly divisible by x + 2.
Solution:
Let the required quadratic polynomial be p(x) = ax2 + bx + c
4. Find the linear polynomials in x, when divided by (x – 3) leaves 6 as remainder and is also exactly
divisible by (x + 3).
Solution:
Let the linear polynomial be f(x) = ax + b
Given p(3) = 6, p(–3) = 0
a(3) + b = 6 ; a(–3) + b = 0
3a + b = 6 ; –3a + b = 0
(1) (2)
Adding (1) and (2)
3a + b = 6
–3a + b = 0
2b = 6 b = 3
Substituting b in (1), we have
3a + 3 = 6 3a = 3
a=1
The required linear polynomial is x + 3.
5. A quadratic polynomial in x leaves remainders 4 and 7 respectively when divided by (x + 1) and (x
– 2). Also it is exactly divisible by (x – 1). Find the quadratic polynomial.
Solution:
Let the quadratic polynomial be p(x) = ax2 + bx + c
Given p(–1) = 4, p(1) = 0 and p(2) = 7
consider p(–1)
a(–1)2 + b(–1) + c = 4
a–b+c=4 (1)
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MATHEMATICS
Now p(1) = 0
and p(2) = 7
a(1)2 + b(1) + c = 0;
a(2)2 + b(2) + c = 7
a +b + c = 0 (2)
4a + 2b + c = 7 (3)
subtracting (2) from (1), we have
a–b+c=4
a+b+c=0
–2b = 4
b = – 2 subtracting (2) from (3), we have
4a + 2b + c = 7
a + b+c=0
3a + b = 7
3a – 2 = 7 (5" b = – 2)
3a = 9
a=3
Substituting a and b in (1), we get c = –1
Hence, the required quadratic polynomial is 3x2 – 2x – 1
6. Find the common factor of the quadratic polynomials 3x2 – x – 10 and 2x2 – x – 6.
Solution:
Consider p(x) = 3x2 – x – 10 and q(x) = 2x2 – x –6
Let (x – k) be the common factor of p(x) and q(x)
p(k) = q(k) = 0
3k2 – k – 10 = 2k2 – k – 6
k2 – 4 = 0
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k2 = 4
k= 4
k=2
The required common factor is (x – 2).
7. Find the remainder when x999 is divided by x2 – 4x + 3
Solution:
Let q(x) and mx + n be the quotient and remainder when x999 is divided by
x2 –4x + 3.
x999 = (x2 – 4x + 3) q(x) + mx + n
If x = 1
199 = (1 – 4 + 3) q(x) + m(1) + n
1 = 0 q(x) + m + n
m + n = 1 (1)
If x = 3
3999 = (32 – 4(3) + 3) q(x) + 3m + n
3999 = 0 q(x) + 3m + n
3m + n = 3999 (2)
subtracting (1) and (2)
3m + n = 3999
m+n=1
2m = 3999 – 1
1 999
m= (3 – 1)
2
Substituting m in (1), we have
1 999
n=1– (3 – 1)
2
1 999 1
=1– 3 +
2 2
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MATHEMATICS
3 1 999
= 3
2 2
3
n= (1 – 3998)
2
1 999 3
The required remainder is (3 – 1) x + (1 – 3998).
2 2
8. Factorize a2(b – c) + b2(c – a) + c2(a – b)
Solution:
Let f(a) = a2(b – c) + b2(c – a) + c2(a – b) —— (1)
put a = b
f(b) = a2(b – c) + b2(c – b) + c2(a – a)
= a2(b – c) – b2(b – c) + 0
f(b) = 0
(a – b) is the factor of f(a).
The given expression is cyclic, so even the factors will be cyclic
Thus the other two factors are (b – c) and (c – a)
a2(b – c) + b2(c – a) + c2(a – b) = k(a – b) + (b – c) + (c – a)
To find the value of a, put a = 0, b = 1, c = –1
0(1+1)+1(–1 –0)+1(0 –1) = k(0 –1) (1+1) (–1 –0)
0 –1 –1 = k(–1) (2) (–1)
–2 = 2k k = –1
The factors of given expression are – (a–b) (b–c) (c–a).
9. Find the remainder when x5is divided by x3 – 4x.
Solution:
Let q(x) be the quotient when x5 is divided by x3 –4x i.e., x(x–2) (x+2)
x5 = (x3 –4x) q(x) + –!x2 +mx + n
put x = 0
0 = 0 q(x) + –!(0) + m(0) + n
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n=0
put x = 2
25 = (8 – 8) q(x) + –!(2)2 + m(2) + n
32 4–! + 2m + n
4–! + 2m = 32 (5"n = 0)
2–! + m = 16 —— (1)
put x = –2
(–2)5 = (–8 + 8) q(x) + –!(–2)2 + m (–2) + n
–32 = 4–! – 2m + n
4–! – 2m = –32 (5"n = 0)
2–! – m = –16 - (2)
Adding (1) and (2)
2–! + m = 16
2–! – m = –16
4–! = 0 –! = 0
Substituting –! in (2), we get
2(0) + m = 16
m=6
The required remainder is 0(x2) + 16x + 0 i.e., 16x
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QUADRATIC EQUATIONS
7 AND INEQUATIONS
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MATHEMATICS
Introduction
If a variable occurs in an equation with positive integral powers and the highest index is two, thenit is called
a Quadratic Equation (in that variable).
A quadratic equation is a second degree polynomial in x equated to zero is a quadratic equation. In other
words, for an equation to be a quadratic, the coefficient of x² should not be zero and the coefficients of any higher
power of x should be 0.
The general form of a quadratic equation is ax2 + bx + c = 0, where a ¹ 0 (and a, b, c are real).
Example (1): x2 – 5x + 6 = 0
Example (2): x2 – x – 6 = 0
Example (3): 2x2 +3x – 2 = 0
Example (4): 2x2 + x – 3 = 0
Roots of the equation:
Just as a first degree equation in x has one value of x satisfying the equation, a quadratic equation in x has
two values of x that satisfy the equation. The values of x that satisfy the equation are called the ROOTS of the
equation. These roots may be real or complex.
The roots of the four quadratic equations given above are:
Equation (1) : x = 2 and x = 3
Equation (2) : x = -2 and x = 3
1
Equation (3) : x= and x = –2
2
3
Equation (4) : x = 1 and x =
2
In general, the roots of a quadratic equation can be found in two ways.
(i) by factoring the expression on the left hand side.
(ii) by using the standard formula.
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All expressions may not be easy to factorize, whereas applying the formula is simple and straightforward.
Finding the roots by factorisation
If the quadratic equation ax2 + bx + c = 0 is written in the form, a (x – ) (x – â) = 0, then the
roots of the equation are and â.
To find the roots of a quadratic equation, we should first express it in the form of
(x ) (x â) = 0, i.e., the left hand side ax2 + bx + c of the quadratic equation
ax2 + bx + c = 0 should be factorized.
For this purpose, we should go through the following steps. We will understand these steps with the help
of the equation x2 – 5x + 6 = 0.
Here a = 1, b = –5 and c = 6
First write down b (the coefficient of x) as the sum of two quantities whose product is equal to ac. In this
case, –5 has to be written as the sum of two quantities whose product is 6. We write –5 as (–3) + (–2), becuase
the product of (–3) and (–2) is equal to 6.
Now, rewrite the equation above.
In this case, the given equation can be written as x2 – 3x – 2x + 6 = 0.
Take the first two terms and rewrite them together after taking out their common factor. Similarly, the third
and the fourth terms should be rewritten after taking out their common factor. In this process, we should ensure
that what is left from the first and the second terms (after removing the common factor) is the same as that left
from the third and fourth terms (after removing their common factor).
In this case, the equation can be rewritten as x(x 3) 2(x 3) = 0; now (x 3) is a
common factor.
In this case, if we take out (x – 3) as the common factor, we can rewrite the given equation as (x – 3) (x –
2) = 0
We know that and â are the roots of the given quadratic equation (x – x) (x – ) = 0.
Hence, the roots of the given equation are 3 and 2.
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MATHEMATICS
For equation (2): x2 – x – 6 = 0, the coefficient of x is –1 which can be rewritten as (–3) + (+2), because
the product of (–3) and 2 is –6, which is equal to ‘ac’
(1 multiplied by –6). Then, we can rewrite the equation as (x – 3) (x + 2) = 0 to get the roots as 3 and –
2.
For equation (3): 2x2 + 3x 2 = 0, the co-efficient of x is 3, which can be rewritten as (+4)
+ (–1) so that their product is –4, which is the value of ‘ac’ (–2 multiplied by 2). Then, we can
1
rewrite the equation as (2x – 1)(x + 2) = 0, obtaining the roots as and –2.
2
For equation (4): 2x2 + x 3 = 0, the coefficient of x is 1, which can be rewritten as (+3) + (2)
because their product is 6, which is equal to ac (2 multiplied by 3). Then we can rewrite the
3
given equation as (x 1)(2x + 3) = 0 to get the roots as 1 and .
2
Finding the roots by using the formula
If the quadratic equation is ax2 + bx + c = 0, then we can use the standard formula given below to find out
the roots of the equation.
b b 2 4ac
x=
2a
Sum and product of roots of a quadratic equation
For the quadratic equation ax2 + bx + c = 0, Let and be the roots, then
b
the sum of the roots ( + ) =
a
c
the product of the roots () =
a
These two rules will be very helpful in solving problems on quadratic equation.
Nature of the roots:
We have already mentioned that the roots of a quadratic equation with real coefficients can be real or
complex. When the roots are real, they can be rational or irrational and, also, they can be equal or unequal.
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The expression b2 – 4ac. Since b2 – 4ac determines the nature of the roots of the quadratic equation, it is
called the “DISCRIMINANT” of the quadratic equation.
A quadratic equation has real roots only if b2 – 4ac ³ 0.
If b2 – 4ac < 0, then the roots of the quadratic equation are complex conjugates.
The following table gives us the clear idea about the nature of the roots of a quadratic equation when a, b
and c are all rational.
Condition Nature of roots
when b2 4ac < 0 the roots are complex conjugates
when b2 4ac = 0 the roots are rational and equal.
when b2 4ac > 0 and a perfect square the roots are rational and unequal.
when b2 4ac > 0 and not a perfect
the roots are irrational and unequal.
square
Note:
1. Whenever the roots of the quadratic equation are irrational, (a, b, c being rational), they are of the
form a + b and a - b , i.e., whenever a + is one root of a quadratic equation, a - is the other root
of the quadratic equation and vice-versa, i.e. if the roots of a quadratic equation are irrational, then
they are conjugate to each other.
2. If the sum of the coefficients of a quadratic equation is zero, then its roots are
c
1 and .
a
c
That is if a + b + c = 0, then the roots of ax2 + bx + c = 0 are 1 and .
a
Signs of the roots:
We can comment on the signs of the roots, i.e. whether the roots are positive or negative, based on the sign
of the sum of the roots and the product of the roots of the quadratic equation. The following table will make it
clear the relationship between the signs of the sum and the product of the roots and the signs of the roots
themselves.
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MATHEMATICS
Sign of product Sign of sum of
Sign of the roots
of the roots the roots
+ ve + ve Both the roots are positive
+ ve ve Both the roots are negative
One root is positive and the other
ve + ve negative. The numerically greater
root is positive
One root is positive and the other
ve ve negative. The numerically greater
root is negative
Constructing a quadratic equation
We can build a quadratic equation in the following cases:
(i) when the roots of the quadratic equation are given.
(ii) when the sum of the roots and the product of the roots of the quadratic equation are given.
Case (i): If the roots of the quadratic equation are and , then its equation can be
written as (x – ) (x – ) = 0 i.e., x2 – x ( + ) + = 0
Case (ii): If p is the sum of the roots of the quadratic equation and q is their product, then the equation can
be written as x2 – px + q = 0.
Constructing a new quadratic equation by changing the roots of a given quadratic equa-
tion:
If we are given a quadratic equation, we can build a new quadratic equation by changing the roots of this
equation in the manner specified to us.
For example, let us take a quadratic equation ax2 + bx + c = 0 and let its roots be and respectively.
Then, we can build new quadratic equations as per the following patterns:
1 1
(i) A quadratic equation whose roots are and , i.e., the roots are reciprocal to the roots of the
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given quadratic equation can be obtained by substituting (1/x) for x in the given equation, which gives us cx2 + bx
+ a = 0, i.e., we get the equation required by interchanging the coefficient of x2 and the constant term.
(ii) A quadratic equation whose roots are (+ k) and (â + k) can be obtained by substituting (x - k) for
x in the given equation.
(iii) A quadratic equation whose roots are (á - k) and (â - k) can be obtained by substituting (x + k) for
x in the given equation.
(iv) A quadratic equation whose roots are (k) and (k) can be obtained by substituting
x
with in the given equation.
k
(v) A quadratic equation whose roots are and can be obtained by substituting
k k
(kx) for x in the given equation.
(vi) A quadratic equation whose roots are (–) and (–) can be obtained by replacing x by
(–x) in the given equation.
Find the roots of the quadratic equation by Graphical method:
First, let us learn how to draw the graph of y = x2
We assume certain real values for x, i.e. we substitute some values for x in y = x2. We can find the
corresponding values of y. We tabulate the values, as shown below.
x 5 4 3 2 1 0 –1 –2 –3 –4 –5
y= x2 25 16 9 4 1 0 1 4 9 16 25
Plotting the points corresponding to the ordered pairs (5, 25), (4, 16), (3, 9),
(2, 4), (1, 1), (0, 0) (–1, 1), (–2, 4), (–3, 9), (–4, 16) and (–5, 25) on the graph paper and joining them
with a smooth curve we obtain the graph of y = x2, as shown below.
We observe the following about the graph of y = x2
1. It is a U shaped graph and it is called a parabola. The arms of the ‘U’ spread outwards.
2. For every value of x (0) we notice that y is always positive. Hence, the graph
lies entirely in the first and second quadrants.
3. When x = 0, y = 0 y = x2 passes through origin
4. The graph is symmetric about the y-axis.
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MATHEMATICS
5. Using the graph of y = x2, we can find the square of any real number as well as the square
root of any non-negative real number
Y
25 y = x2
20
15
10
5
X' X
–5 –4 –3 –2 –1 0 1 2 3 4 5
Y'
(a) for any given x value, the corresponding y value on the graph is its square
and
(b) for any given y ( 0) value, the corresponding x value on the graph is its square root.
6. The graph of y = kx2, when k > 0 lies entirely in Q1 and Q2 and when k < 0 the graph lies entirely in
Q3 and Q4
Y Y
0 X
X
0
y = kx2, k > 0 y = kx2, k < 0
The method of solving the quadratic equation of the form px2 + qx + r = 0 whose roots are real is shown
in the following example.
Solve x2 – 5x + 6 = 0 using the graphical method
Solution:
Let y = x2 – 5x + 6
Prepare the following table by assuming different values for x
x 0 1 2 3 –1 –2 4 5
x2 0 1 4 9 1 4 16 25
5x 0 5 10 15 –5 10 20 25
y= x2 – 5x + 6 6 2 0 0 12 20 2 6
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Plot the points (0, 6), (1, 2), (2, 0), (3, 0), (–1, 12), (–2, 20), (4, 2) and (5, 6) on the graph and join the
points with a smooth curve, as shown below.
Y
12
10
8
6
4
2
X' X
–3 –2 –1 0 1 2 3 4 5
–2
–4
–6
Y'
Here we notice that the given graph (parabola) intersects the x-axis at (2, 0) and (3, 0)
The roots of the given quadratic equation x2 –5x + 6 = 0 are x = 2 and x = 3.
The roots of the given equation are the x-coordinates of the points of intersection of the curve with x-
axis.
Note:
1. If the graph meets the x-axis at two distinct points, then the roots of the given equation are real and
distinct
X
0 (x2, 0)
(x1, 0)
2. If the graph touches the x-axis at only one point, then the roots of the quadratic equation are real and
equal
0 X
(x, 0)
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MATHEMATICS
3. If the graph does not meet the x-axis, then the roots of the quadratic equation are not real, i.e. they
are complex
0 X
Second method:
We can also solve the quadratic equation px2 + qx + r = 0 by considering the following equations y = px2
––––– (1) and y = – qx – r –––––– (2)
Clearly, y = px2 is a parabola and y = – qx – r is a straight line.
Step (i) Draw the graph of y = px2 and y = – qx – r on the same graph paper.
Step (ii) Draw perpendiculars from the points of intersection of parabola and the line onto the x-axis. Let
the points of intersection on the x-axis be (x1, 0) and (x2, 0).
Step (iii) The x-coordinates of the points in step (ii), i.e. x1 and x2 are the two distinct roots of px2 + qx =
r=0
Example:
Solve 2x2 – x – 3 = 0
Solution:
We know the roots of 2x2 – x – 3 = 0 are the x coordinates of the points of intersection of the
parabola y = 2x2 and the line y = x + 3
(1) y = 2x2 (2) y=x+3
X 0 1 2 –1 –2 x 0 1 2 –1 –2 –3
y = 2x2 0 2 8 2 8 y=x+3 3 4 5 2 1 0
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Draw the graph of y = 2x2 and y = x + 3
Y
y = x2
8
y=x+3
6
4
2
X' X
–3 –2 –1 0 1 1.5 2 3
–2
Y'
Clearly, the perpendiculars drawn from the points of intersection of parabola and the line meet the x-axis
3
at , 0 and (–1, 0)
2
3
The roots of the given quadratic equation 2x2 – x – 3 = 0 are and –4
2
Note:
(i) If the line meets the parabola at two points, then the roots of the quadratic equation are real and
distinct
(ii) If the line touches the parabola at only one point, then the quadratic equation has real and equal roots
(iii) If the line does not meet the parabola, i.e. when the line and the parabola have no points in common,
then the quadratic equation has no real roots. In this case, the roots of the quadratic equation are
imaginary
Y Y
X X
0 one root 0
Equations of higher degree:
The index of the highest power of x in the equation is called the degree of the equation. For example, if the
highest power of x in the equation is x3, then the degree of the equation is 3. An equation whose degree is 3 is
called a cubic equation. A cubic equation will have three roots.
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MATHEMATICS
An equation whose degree is n will have n roots.
Maximum or minimum value of a quadratic expression:
The quadratic expression ax2 + bx + c takes different values as x takes different values.
For all the values of x, as x varies from to +, (i.e. when x is real), the quadratic expression
ax2 + bx + c
(i) has a minimum value if a > 0 (i.e. a is positive). The minimum value of the
quadratic expression is
4ac b and it occurs at x = b .
2
4a 2a
(ii) has a maximum value if a < 0 (i.e. a is negative). The maximum value of the
quadratic expression is
4ac b and it occurs at x = b .
2
4a 2a
Examples
1. Find the roots of the equation x² + 3x 4 = 0.
Solution:
x² + 3x 4 = 0 x² x + 4x 4 = 0
x(x 1) + 4(x 1) = 0 (x + 4) (x 1) = 0
x = 4 or x = 1
2. Find the roots of the equation 4x² 13x + 10 = 0
Solution:
4x² 13x + 10 = 0
4x² 8x 5x + 10 = 0
4x(x 2) 5(x 2) = 0 (4x 5) (x 2) = 0
5
x= or x = 2.
4
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3. Find the roots of the equation 26x² 43x + 15 = 0
Solution:
We have to write 43 as the sum of two parts whose product should be equal to (26)
(15).
26 15 = 13 30 and 13 + 30 = 43
26x² 43x + 15 = 0
26x² 13x 30x + 15 = 0
(13x 15) (2x 1) = 0
15 1
x= or x =
13 2
We can also find the roots of the equation by using the formula.
b b 2 4ac
x=
2a
=
43 432 1560 =
43 1849 1560
52 52
43 289 43 17
= x=
52 52
43 17 43 17 60 26
x= or = or
52 52 52 52
15 1
x= or
13 2
4. Discuss the nature of the roots of the equation 4x² - 2x + 1 = 0.
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MATHEMATICS
Solution:
Discriminant = (2)² 4(4)(1) = 4 16 = 12 < 0
Since the discriminant is negative, the roots are imaginary.
3
5. If the sum of the roots of the equation kx² 3x + 9 = 0 is , then find the product of
11
the roots of that equation.
Solution:
3 3
Sum of roots of the equation = = (given)
k 11
k = 11
9
In the given equation, product of roots =
k
9
As k = 11, product of roots =
11
6. Form the quadratic equation whose roots are 2 and 7.
Solution:
Sum of the roots = 2 + 7 = 9
Product of roots = 2 7 = 14
We know that if p is the sum of the roots and q the product of the roots of a quadratic equation, then its
equation is x² - px + q = 0
Hence the required equation is x² - 9x + 14 = 0
7. Form a quadratic equation with rational coefficients one of whose roots is
3+ 5.
Solution:
If (3 + 5 ) is one root, then the other root is (3 5 )
Sum of the roots = 6
Product of the roots = 4
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Thus, the required equation is x² 6x + 4 = 0
8. A person can buy 15 books less for Rs.900 when the price of each book goes up by Rs.3. Find the
original price and the number of copies he could buy at the initial price.
Solution:
900
Let the number of books bought initially for Rs.900 be ‘x’. The original price of book was .
x
Now the price of the book is increased by Rs.3
900
i.e., the new cost is Rs. +3
x
And the number of books bought is reduced by 15 i.e., (x - 15)
Since the total amount spent is still Rs.900, the product of the price and the number of books are still 900.
900
x 3 (x 15) = 900
(900 + 3x) (x 15) = 900x
3x² + 855x 13500 = 900x
3x² 45x 13500 = 0 x² 15x 4500 = 0
x² 75x + 60x 4500 = 0
x(x 75) + 60(x 75) = 0
(x 75) (x + 60) = 0 x = 75 or 60
Since x cannot be negative, x = 75
900
Thus, the original price of the book = = Rs.12
75
9. If and are the roots of the equation x² 6x + 8 = 0, then find the values of
(i) ² + ²
1 1
(ii) +
(iii) ( > )
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MATHEMATICS
Solution:
From the given equation, we get + = 6 and = 8
(i) ² + ² = ( + )² 2 = (6)² 2(8) = 20
1 1 6 3
(ii) + = = =
8 4
(iii) ( )² = ( + )² 4
( ) = 2 4 = 6 2 4(8) ( ) = 2
– = 2, (5" > )
10. Find the value of x given 3x+1 + 32x+1 = 270
Solution:
3x+1 + 32x+1 = 270 3.3x + 32x.3 = 270
3x + 32x = 90
Substituting 3x = a, we get,
a + a² = 90 a² + a 90 = 0
a² + 10a 9a 90 = 0
(a + 10) (a – 9) = 0 a = 9 or a = 10
If 3x = 9, then x = 2.
If 3x = 10, which is not possible.
x=2
11. Solve |x|2 7 |x| | 12 = 0.
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Solution:
Given equation is |x|2 7 |x| + 12 = 0
(|x| 3) (|x| 4) = 0
|x| = 3 or |x| = 4
x = 3 or x = 4.
12. Solve |x|2 + 7 |x| + 10 = 0.
Solution:
Given equation is |x|2 +7 |x| + 10 = 0
(|x| + 2) (|x| + 5) = 0
|x| = 2 or |x| = 5
But, absolute value of any number can never be negative.
No roots are possible for the given equation.
Quadratic Inequations:
Consider the quadratic equation ax2 + bx + c = 0, (a 0) where a, b and c are real numbers.
The quadratic inequalities related to ax2 + bx + c = 0 are ax2 + bx + c < 0 and
ax2 + bx + c > 0.
Assume that a 0.
The following cases arise:
Case (i):
If b2 4ac 0, then the equation ax2 + bx + c = 0 has real and unequal roots.
Let and ( ) be the roots.
Then,
ax2 + bx + c = a (x ) (x )
–
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MATHEMATICS
(a) If x , then (x ) 0 and (x ) 0
ax2 + bx + c 0
(b) If x , then (x ) 0 and (x ) 0
ax2 + bx + c 0
(c) If x , then x 0 and x 0
ax2 + bx + c 0
Case (ii):
If b2 4ac = 0, then ax2 + bx + c = 0 has real and equal roots.
Let x1 be the equal root.
ax2 + bx + c = a(x x1) (x x1)
– x1
(a) If x x1. Then x x1 0
ax2 + bx + c 0
(b) If x x1, then x x1 0
ax2 + bx + c 0
Case (iii):
If b2 4ac 0, then ax2 + bx + c = 0 has imaginary roots.
In this case, ax2 + bx + c 0, x R.
The above concept can be summarised as
(i) If < x < , then (x – ) (x – ) < 0 and vice-versa.
(ii) If x < and x > ( < ), then (x – ) (x – ) > 0 and vice-versa.
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Examples:
Example 1:
Solve the inequation x2 + x - 6 < 0
Solution:
Given inequation is x2 + x 6 0
(x + 3) (x 2) 0
(x + 3)< 0, (x – 2) > 0 or (x + 3) > 0, (x – 2) < 0
x < –3, x > 2 (case I) (or)
x > –3, x < 2 (case II)
Case I:
x < –3 and x > 2
–5 –4 –3 –2 –1 0 1 2 3 4 5
x < –3 x>2
There exists no value of x so that x < –3 and x > 2 (as there is no overlap of the regions), hence no value
of x in this case satisfies the given inequation,
Case II:
x > –3 and x < 2
–5 –4 –3 –2 –1 0 1 2 3 4 5
x > –3
x<2
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All the points in the overlapping region, i.e. –3 < x < 2, satisfy the inequation.
Hence, the solution set of the inequation
x2 + x 6 0 is, {x / 3 x 2} or (3, 2).
Example 2:
Solve for x: x2 4x 3 0
Solution:
Given inequation is x2 4x 3 0
(x 1) (x 3) 0
x – 1 e” 0; x – 3 e” 0 or x – 1 d” 0; x – 3 d” 0
x e” 1; x e” 3 (case I) (or)
x d” 1; x d” 3 (case II)
Case I:
x ³ 1 and x ³ 3
–3 –2 –1 0 1 2 3 4
x1
All the points in the overlapping region, i.e. x 3, satisfy the given inequation.
Case II:
x 1 and x 3
–4 –3 –2 –1 0 1 2 3
x1
x3
All the points in the overlapping region, i.e. x 1, satisfy the given inequation
Hence, the solution per the given inequation is
x (–, 1] [3, )
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Example 3:
Solve x2 + 6x + 13 0
Solution:
Given inequation is x2 + 6x + 13 0
Here, factorisation is not possible.
Rewriting the given inequation by identifying the hidden square, we get
(x2 + 6x + 9) + 4 0 (x + 3)2 + 4 > 0.
We know that (x + 3)2 0 x R,
(x + 3)2 + 4 4 0 x R,
The required solution is the set of all real numbers, i.e. (, ).
Example 4:
x 2 5x 3
Solve <x
x2
Solution:
x 2 5x 3
x
x2
x 2 5x 3 x 2 5x 3 x 2 2 x
x0 0
x2 x2
3x 3
x 2 0
x 1
x 2 0 —— (1)
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MATHEMATICS
To solve (1), it is sufficient to solve (x + 1) (x + 2) < 0.
We know that (x – ) (x – ) < 0 < x < ( < )
–2 x –1
Thus, the required solution is –2 x –1
Example 5:
1 2
Solve x 1 1 2x
Solution:
1 2
x 1 1 2x
1 2
0
x 1 1 2x
1 2x 2x 2
( x 1) (1 2 x ) < 0
1
( x 1) (1 2 x ) < 0
1
( x 1) (1 2 x ) > 0 —— (1)
(1) holds good if (x – 1) (1 – 2x) > 0
We know that (x – ) (x – ) > 0 ( < ) x < , x >
1
The solution of the given inequation is x and x > 1,
2
1
i.e. x (1, ).
2
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8 PROGRESSIONS
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MATHEMATICS
Sequences of numbers which follow specific patterns are called progression. Depending on the pattern,
the progressions are classified as follows.
(i) Arithmetic Progression
(ii) Geometric Progression and
(iii) Harmonic Progression
ARITHMETIC PROGRESSION (A.P.)
Numbers (or terms) are said to be in arithmetic progression when each one, except the first, is obtained by
adding a constant to the previous number (or term).
An arithmetic progression can be represented by a, a + d, a + 2d, …, [a + (n – 1)d]. Here, d is added to
any term to get the next term of the progression. The term a is the first term of the progression, n is the number of
terms in the progression and d is the common difference. The nth term is normally represented by Tn and the sum
to n terms of an A.P. is denoted by Sn
nth term = Tn = a + (n – 1)d
n
Sum to n terms = Sn = [2a + (n – 1)d]
2
The sum to n terms of an A.P. can also be written in a different manner. That is,
n n
sum of n terms = [2a + (n – 1)d] = [a + {a + (n – 1)d}]
2 2
But, when there are n terms in an A.P., a is the first term and {a + (n – 1)d} is the last term. Hence,
n
Sn = [first term + last term]
2
The average of all the terms in an A.P. is called the arithmetic mean (A.M.) of the A.P. Since the average of
a certain numbers is equal to the {sum of all the number/number of numbers}.
1 n
A.M. of n terms in A.P. = Snn =
n 2
(First Term Last Term)
(First Term + Last Term) =
2
i.e. The A.M. of an A.P. is the average of the first and the last terms of the A.P.
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The A.M. of an A.P. can also be obtained by considering any two terms which are EQUIDISTANT from
the two ends of the A.P. and taking their average, i.e.
(a) the average of the second term from the beginning and the second term from the end is equal to the
A.M. of the A.P.
(b) the average of the third term from the beginning and the third term from the end is also equal to the
A.M. of the A.P. and so on.
In general, the average of the kth term from the beginning and the kth term from the end is equal to the A.M.
of the A.P.
If the A.M. of an A.P. is known, the sum to n terms of the series (Sn) can be expressed as Sn = n (A.M.)
In particular, if three numbers are in arithmetic progression, then the middle number is the A.M. i.e. if a, b
ac
and c are in A.P., then b is the A.M. of the three terms and b .
2
ab
If a and b are any two numbers, then their A.M. = .
2
Note:
(i) If three numbers are in A.P., we can take the three terms to be (a – d), a and
(a + d).
(ii) If four numbers are in A.P., we can take the four terms to be (a – 3d), (a – d), (a + d) and (a + 3d).
The common difference in this case is 2d and not d.
(iii) If five numbers are in A.P., we can take the five terms to be (a – 2d), (a – d),
a, (a + d) and (a + 2d).
Inserting arithmetic mean between two numbers:
When n arithmetic means a1, a2, .........., an are inserted between a and b, then a, a1, a2, ........, an, b are in
A.P. t1 = a and tn+2 = b of A.P.
The common difference of the A.P. can be obtained as follows:
Given that, n arithmetic means are there between a and b.
a = t1 and b = tn + 2
Let d be the common difference.
b = t1 + (n + 1) d
b = a + (n + 1)d
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MATHEMATICS
(b a )
d=
(n 1)
Some important results:
The sum to n terms of the following series are quite useful and, hence, should be remembered by students.
n n(n 1)
(i) Sum of the first n natural numbers = i
1 2
n n(n 1)(2n 1)
2
(ii) Sum of squares of the first n natural numbers = i
1 6
(iii) Sum of cubes of first n natural numbers
2 2
n
3 n(n 1) n 2 (n 1) 2 n
= i 2 i
1 4 1
Example 1:
Find the 14th term of an A.P. whose first term is 3 and the common difference is 2.
Solution:
The nth term of an A.P. is given by tn = a + (n 1)d, where a is the first term and
d is the common difference.
t14 = 3 + (14 1) 2 = 29
Example 2:
Find the first term and the common difference of an A.P. if the 3rd term is 6 and the 17th term is 34.
Solution:
If a is the first term and the common difference d, then we have
a + 2d = 6 ——— (1)
a + 16d = 34 —— (2)
On subtracting equation (1) from equation (2), we get
14d = 28 d = 2
Substituting the value of d in equation (1), we get a = 2
a = 2 and d = 2
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Example 3:
Find the sum of the first 22 terms of an A.P. whose first term is 4 and the common difference is 4/3.
Solution:
4
Given that, a = 4 and d = .
3
n
We have Sn = [2a + (n – 1)d]
2
22 4
S22 = 2 4 22 1 11 8 28 396
2 3
Example 4:
Divide 124 into four parts in such a away that they are in A.P. and the product of the first and the 4th
part is 128 less than the product of the 2nd and the 3rd parts.
Solution:
Let the four parts be (a 3d), (a d), (a + d) and (a + 3d). The sum of these four parts is
124,
i.e. 4a = 124 a = 31
(a 3d) (a + 3d) = (a d) (a + d) 128
a² 9d² = a² d² 128
8d² = 128 d = 4
As a = 31, taking d = 4, the four parts are 19, 27, 35 and 43.
Note:
If d is taken as –4, then the same four numbers are obtained, but in decreasing order.
Example 5:
Find the three terms in A.P., whose sum is 36 and product is 960.
Solution:
Let the three terms of an A.P. be (a d), a and (a + d).
Sum of these terms is 3a.
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MATHEMATICS
3a = 36 a = 12
Product of these three terms is
(a + d) a (a d) = 960 (12 + d) (12 d) = 80
144 d² = 80 d = 8
Taking d = 8, we get the terms as 4, 12 and 20.
Note:
If d is taken as -8, then the same numbers are obtained, but in decreasing order.
Geometric Progression (G.P.):
Numbers are said to be in geometric progression when the ratio of any quantity to the number that follows
it is the same. In other words, any term of a G.P. (except the first one) can be obtained by multiplying the previous
term by a constant factor.
The constant factor is called the common ratio and is normally represented by r. The first term of a G.P. is
generally denoted by a.
A geometric progression can be represented by a, ar, ar2, ..... where a is the first term and r is the common
ratio of the G.P. nth term of the G.P. is arn – 1 i.e. tn = arn–1
Sum to n terms = Sn =
a 1 r n
=
a r n 1
=
r ar n 1 a
1 r r 1 r 1
The sum to n terms of a geometric progression can also be written as
r Last term First term
Sn =
r 1
Note:
If n terms a1, a2, a3, ......... an are in G.P., then the geometric mean (G.M.) of these n terms is given by =
n a 1a 2a 3 ........... a n
If three terms are in geometric progression, then the middle term is the geometric mean of the G.P., i.e. if a,
b and c are in G.P., then b is the geometric mean of the three terms.
If there are two terms a and b, their geometric mean is given by G.M. = ab . We see that the 3 terms
a, ab , b are in G.P..
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For any two positive numbers a and b, their arithmetic mean is always greater than or equal to their
geometric mean, i.e. for any two positive numbers
ab
a and b, e” ab . The equality holds if and only if a = b.
2
Note:
When there are three terms in geometric progression, we can take the three terms to be a/r, a and ar.
Infinite geometric progression:
If –1 < r < 1 (or | r | < 1), then the sum of a geometric progression does not increase infinitely but
“converges” to a particular value, no matter how many terms of the G.P. we take. The sum of an infinite geometric
progression is represented by S and is given by the formula,
a
S = , if r < 1.
1 r2
Example 6:
Find the 7th term of the G.P. whose first term is 6 and common ratio is 2/3.
Solution:
2
Given that, t1 = 6 and r =
3
We have tn = a . rn1
6
2 (6) (64) 128
t7 = (6) = =
3
729 243
Example 7:
Find the common ratio of the G.P. whose first and last terms are 25 and 1/625 respectively and the sum of
the G.P. is 19531/625.
Solution:
first term r last term
We know that the sum of a G.P is 1r
19531 25 r / 625
r = 1/5
625 1 r
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MATHEMATICS
Example 8:
Find three numbers of a G.P. whose sum is 26 and product is 216.
Solution:
Let the three numbers be a/r, a and ar.
Given that,
a/r . a . ar = 216;
a3 = 216; a = 6
a/r + a + ar = 26
6 + 6r + 6r² = 26r
6r² 20r + 6 = 0
6r² 18r 2r + 6 = 0
6r (r 3) 2(r 3) = 0
r = 1/3 (or) r = 3
Hence the three numbers are 2, 6 and 18 (or) 18, 6 and 2
Example 9:
If | x | < 1, then find the sum of the series 2 + 4x + 6x² + 8x3 + . . .
Solution:
Let S = 2 + 4x + 6x² + 8x3 + ——— (1)
xS = 2x + 4x2 + 6x3 + …..(2)
(1) – (2) gives
S (1 x) = 2 + 2x + 2x² + 2x3 + . . .
= 2 (1 + x + x² + . . .)
1 + x + x² + . . . is an infinite G.P with a = 1, r = x and |r| = |x| < 1
Sum of the series = 1/1x
S (1 x) = 2/(1 x)
S = 2/(1 x)²
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Example 10:
Find the sum of the series 1, 2/5, 4/25, 8/125, . . . . . . .
Solution:
2 2
Given that, a = 1, r = and |r| = <1
5 5
a 1
S = 1 r 1 2 = 5/3
5
Note:
When n geometric means are there between a and b, the common ratio of the G.P. can be derived as
follows.
Given that, n geometric means are there between a and b.
a = t1 and b = tn + 2
Let ‘r’ be the common ratio
b = (t1) (rn+1) b = a rn+1
b
rn + 1 =
a
( n 1) b
r=
a
Harmonic Progression (H.P.):
A progression is said to be a harmonic progression if the reciprocal of the terms in the progression form an
arithmetic progression.
For example, consider the series
1 1 1 1
, , , ,...
2 5 8 11
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MATHEMATICS
The progression formed by taking reciprocals of terms of the above series is 2, 5, 8, 11,. . . . Clearly, these
terms form an A.P. whose common difference is 3.
Hence, the given progression is a harmonic progression.
nth term of an H.P:
We know that if a, a + d, a + 2d,….are in A.P., then the nth term of this A.P. is
a + (n – 1) d. Its reciprocal is
1
a n 1d
1 1 1
So, nth term of an H.P. whose first two terms are and a d is a n 1d
a
Note:
There is no concise general formula for the sum to n terms of an H.P.
Example 11:
3 3 3
Find the 10th term of the H.P. , 1, , , ………….
2 4 5
Solution:
3 3 3
The given H.P. is , 1, , ,........
2 4 5
2 4 5
The corresponding A.P. is , 1, , ,.......
3 3 3
2 2 1
Here a = ;d=1– =
3 3 3
2 1 11
T10 of the corresponding A.P. is a + (10 – 1)d = (9)
3 3 3
3
Hence required term in H.P. is
11
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Harmonic Mean (H.M.):
If three terms are in H.P., then the middle term is the H.M. of other two terms.
The harmonic mean of two terms a and b is given by
2ab
H.M.
ab
Relation between A.M., H.M. and G.M. of two numbers
Let x and y be two numbers
xy 2xy
A.M. = , G.M. = xy and H.M. =
2 xy
(A.M.) (H.M.) = (G.M)2
Inserting n harmonic means between two numbers:
To insert n H.M.’s between two numbers, we first take the corresponding arithmetic series and insert n
arithmetic means, and next, we find the corresponding harmonic series.
This is illustrated by the example below:
Example 12:
1 1
Insert three harmonic means between and
12 20
Solution:
After inserting the harmonic means
let the harmonic progression be
1 1 1 1 1
, , , ,
a a d a 2d a 3d a 4d
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MATHEMATICS
1 1 1 1
As = and
a 12 a 4d 20 a = 12 and d = 2
=
1 1 1
The required harmonic means are , and
14 16 18
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9 STATEMENTS
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MATHEMATICS
Introduction
A sentence with can be judged either true or false but not both is called a statement. Statements are
denoted by lower case letters p, q, r etc.
Examples:
1. p : 2 is a prime number. This statement is true.
2. q : 2 + 3 = 6. This statement is false
Truth value:
The truthness or falsity of a statement is called its truth value.
Examples:
1. The truth value of the statement p : The sun rises in the east, is True.
2.The truth value of the statement q : All odd numbers are prime, is False.
Truthness of a statement is denoted by T, while its falsity is denoted by F.
Negation of a statement:
The denial of a statement is called its negation. Negation of a statement p is denoted by ~ p and read
as “not p” or “negation p”.
Truth table:
p p
T F
F T
Examples:
1. p:2+4=6
p:2+46
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2. p : 3 is a factor of 10
p : 3 is not a factor of 10
3. p : Charminar is in Delhi
p : Charminar is not in Delhi
Compound statement:
A statement obtained by combining two or more simple statements using connectives is called a compound
statement.
Examples:
Consider the two statements.
p : 2 is a prime number and
q : 2 is an even number.
Some compound statements that can be formed by using the statements p and q are:
(i) 2 is a prime number and 2 is an even number.
(ii) 2 is a prime number or 2 is an even number.
(iii) 2 is neither a prime number nor an even number.
Let us look at some basic compound statements.
1. Conjunction:
If p and q are two simple statements, the compound statement “p and q” is called the
conjunction of p and q. It is denoted by p q.
Truth table:
p q pq
T T T
T F F
F T F
F F F
We observe that p q is true only when both p and q are true.
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MATHEMATICS
Examples:
(i) Let p : 4 is a perfect square
q : 2 is an odd number.
p q : 4 is a perfect square and 2 is an odd number.
As p is T and q is F, the truth value of p q is F.
(ii) Let p : 3 2
q : 2 is an irrational number.
Then,
p q : 3 2 and is an irrational number.
The truth value of pq is true as both p and q are true.
2. Disjunction:
If p and q are two simple statements, then the compound statement “p or q” is called the disjunction of p
and q. It is denoted by p q.
Truth table:
p q pq
T T T
T F T
F T T
F F F
We observe that, p q is false only when both p and q are false.
Examples:
(i) Let p : The set of even primes is an empty set.
q : 1 is a factor of every natural number.
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p q : The set of even primes is an empty set or 1 is a factor of every natural
number.
As p is false and q is true, truth value of p q is T.
(ii) Let p : 5 is a factor of 18
q : 12 divides 6.
Then, p q : 5 is factor of 18 or 12 divides 6.
The truth value of p q is false.
3. Implication or Conditional:
If p and q are two statements, the compound statement “if p then q”, is called a conditional statement. It
is denoted by p q.
The statement p is called the hypothesis (or given) and the statement q is called the conclusion (or result).
Truth table:
p q pq
T T T
T F F
F T T
F F T
We observe that, a true statement cannot imply a false statement.
Examples:
(i) Let p : Every set is a subset of itself.
q:3+5=8
p q : If every set is a subset of itself, then 3 + 5 = 8.
As p is true and q is true, the truth value of p q is true.
(ii) Let p : ABC is a right triangle it A = 100°
q : A + B + C = 180°
p q : If ABC is a right triangle if A = 100°, then A + B + C = 180°
As p is false and q is true, the truth value of p q is true.
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MATHEMATICS
4. Bi - conditional or Bi - implication:
If p and q are two statements, then the compound statement “p if and only if q” is called the bi-conditional
of p and q. It is denoted by p q.
Truth table:
p q pq
T T T
T F F
F T F
F F T
We observe that, p q is true if both p and q have the same truth values.
Examples:
1. Let p : 2 3 = 6
q : 2 + 8 = 10
p q : 2 3 = 6 if and only if 2 + 8 = 10.
Since both p and q are true, the truth value of p q is T.
2. Let p : Every triangle is equilateral
q : Charminar is in Hyderabad
p q : Every triangle is equilateral if and only if Charminar is in Hyderabad.
As p is false and q is true the truth value of p q is F.
Converse, inverse and contrapositive of a conditional :
Let p q or if p then q be a conditional,
(i) If q then p i.e., q p, is called the converse of p q.
(ii) If not p then not q i.e., p q, is called the inverse of p q
(iii) If not q then not p i.e., q p is called the contrapositive of p q.
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Truth table:
Conditional Converse Inverse Contrapositive
p q
p q qp p q q p
T T T T T T
T F F T T F
F T T F F T
F F T T T T
Examples:
1. Write the converse, inverse and contrapositive of the conditional “If x is odd then x2 is odd”.
Solution:
Conditional: If x is odd, then x2 is odd.
Converse : If x2 is odd, then x is odd.
Inverse : If x is not odd, then x2 is not odd.
Contrapositive : If x2 is not odd, then x is not odd.
2. Write the converse, and the contrapositive of the conditional, “If ABC is a triangle, then ÐA + ÐB +
ÐC = 180°”.
Solution:
Conditional : If ABC is a triangle then A +B + C = 180°.
Converse : If A + B + C = 180° then ABC is a triangle.
Contrapositive : If A + B + C 180° then ABC is not a triangle.
Let us now look at the truth tables of some compound statements:
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MATHEMATICS
Examples:
1. The truth table of p ~ q :
p q q pq
T T F T
T F T T
F T F F
F F T T
2. The truth table of ~ p (p q)
p q p pq p (p q)
T T F T T
T F F F F
F T T F T
F F T F T
3. Write the truth table of ~ p p q
Solution:
p q p pq p p q
T T F T T
T F F T T
F T T T T
F F T F F
Tautology:
A compound statement which always takes True as its truth value is called a tautology.
Examples:
1. The truth table of p ~ p is
p p pp
T F T
F T T
We observe that p ~p takes T as its truth value always. So, p ~ p is a tautology.
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2. Show that the compound statement p p q is a tautology.
Solution:
Truth table:
p q pq p (p q)
T T T T
T F T T
F T T T
F F F T
We observe that p pq is always True
Hence, p p q is a tautology.
Contradiction:
A compound statement which always takes False as its truth value is called a contradiction.
For example:
1. The truth table of p ~p is
p p p p
T F F
F T F
We observe that p ~p takes F as its truth value always. So p ~p is a contradiction.
2. Show that the compound statement is a contradiction.
Solution:
Truth table:
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MATHEMATICS
p q p q p p q q p p q q
T T F F T F F
T F F T T F F
F T T F T F F
F F T T T F F
We observe that p p q q is always False.
Hence, p p q q is a contradiction.
Contingency:
A compound statement which is neither tautology nor contradiction is called a contingency.
Example: (p q) <“ p
Truth table:
p q p pq p q p
T T F T F
T F F T F
F T T T T
F F T F T
Logically equivalent statements:
Two statements r and s are said to be logically equivalent, if the last column of their truth tables are
identical.
(OR)
Two statements r and s are said to be logically equivalent if r s is a tautology. Generally, r and s will
be compound statements.
If the statements r and s are logically equivalent, then we denote this as r a” s.
Note that r s is always true only if both r and s have same truth values.
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Examples:
1. Show that p q q p
Solution:
p q pq qp
T T T T
T F F F
F T F F
F F F F
We observe that p q and q p have the same truth values. Hence, p q q p.
2. Show that p q p q
Solution:
Truth table
p q p pq p q
T T F T T
T F F F F
F T T T T
F F T T T
We observe that p q and p q have the same truth values.
Hence, p q p q
Laws of algebra of statements:
Some logical equivalences are listed under the following laws :
1. Commutative Laws:
(a) pqqp
(b) pqqp
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MATHEMATICS
2. Associative laws:
(a) (p q) r p (q r)
(b) (p q) r p (q r)
3. Distributive Laws:
(a) p (q r) (p q) (p r)
(b) p (q r) (p q) (p r)
4. Idempotent Laws:
(a) p p p
(b) p p p
5. De Morgan’s laws:
(a) (p q) (p) (q)
(b) (p q) (p) (q)
6. Identity Laws:
(a) p f a” p, p t a” t.
(b) p f a” f, p t a” p.
7. Complement Laws:
(a) p(p) a” t (b) p (p) a” f (c) (p) p
(d) t a” f (e) f a” t
List of equivalences based on implications:
(i) pq p q
(ii) (p q) pq
(iii) pq q p
(i.e., a conditional and its contrapositive are logically equivalent)
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(iv) q p q p)
(i.e., converse and inverse of a conditional are logically equivalent)
(v) p q (p q) (q p)
Open sentence:
A sentence involving one or more variables is called an open sentence, if it becomes TRUE or FALSE
when the variables are replaced by some specific values from the given set. The set from which the values of a
variable can be considered is called the replacement set or domain of the variable.
Examples:
1. x + 2 = 9 is an open sentence.
For x = 7, it becomes True and for other real values of x it becomes False.
2. x2 + 1 0 is an open sentence.
For all real values of x it is True.
Quantifiers:
A quantifier is a word or phrase which quantifies a variable in the given open sentence.
There are two types of quantifiers.
(a) Universal quantifier.
(b) Existential quantifier.
Universal quantifier:
The quantifiers like “for all”, “for every”, “for each” are called universal quantifiers. A
universal quantifier is denoted by ‘ ’.
Examples:
1. Consider the open sentence, x 0 .
This is true for all x R. So, we write , x R.
2. Consider the sentence,
1 + 2 + 3 + . . . . . + n = , n N.
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MATHEMATICS
Existential quantifier:
The quantifiers like “for some”, “not all”, “there is/exists at least one” are called existential quantifiers.
An existential quantifier is denoted by ‘$’.
Examples:
1. Not all prime numbers are odd.
2. x R such that x + 4 = 11
Negation of statements involving quantifiers:
1. p : All odd numbers are prime.
~p : Not all odd numbers are prime
(or)
Some odd numbers are not prime.
(or)
There is an odd number which is not prime.
2. p : All questions are difficult.
~ p : Not all questions are difficult.
(or)
Some questions are not difficult.
(or)
There is at least one question which is not difficult.
3. p : All birds can fly
~ p : Not all birds can fly
(or)
There are some birds which cannot fly.
(or)
There is at least one bird which cannot fly.
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Methods of Proof:
Statements in mathematics are usually examined for their validity. The various steps involved in the process
is referred as proof.
There are two important types of proofs
1. Direct Proof:
In this method, we begin with the hypothesis and end up with the desired result through a logical se-
quence of steps.
Example:
If x is odd, then x2 is odd.
Given, x is an odd number
Conclusion: x2 is an odd number
Proof:
x = 2k + 1 for some k Z
x2 = (2k + 1)2
= 4k2 + 4k + 1
= 2(2k2 + 2k) + 1
= 2m + 1, where m = 2k2 + 2kz (5" k z 2k2 + 2kz)
x2 = 2m + 1, m Z
x2 is an odd number.
Hence, if x is an odd number, then x2 is an odd number.
2. Indirect proof:
In this method, we proceed by assuming that the conclusion is false. Then we arrive at a contradiction. This
implies that the desired result must be true.
Example:
If a + b = 0, then (a + b)2 = 0. where a, b z – {0}
Given: a + b = 0
Conclusion: (a + b)2 = 0
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MATHEMATICS
Proof:
Let us assume that (a + b)2 0
(a + b)2 0 (As (a + b)2 cannot be –ve)
a + b 0 or a + b 0
which is a contradiction to the hypothesis i.e., a + b=0
Our assumption (a + b)2 0 is flase.
Hence, if a + b = 0, then (a + b)2 = 0.
Methods of Disproof:
To disprove a given statement there are two methods.
1. Counter example method:
In this method, we look for a counter example which disproves the given statement.
Examples:
(i) Every odd number is a prime number.
This statement is false, as 9 is an odd number but it is not a prime.
(ii) x2 – x – 6 = 0 for all real values of x
This statement is false, as for x = 2,
x2 – x – 6 = (2)2 – 2 – 6 = – 4 0
x = 2 is a counter example here.
2. Method of contradiction:
In this method, we assume that the given statement is true. Then we arrive at a contradiction. This implies
that the given statement is false.
Example:
Disprove the statement, “There can be two right angles in a triangle”.
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Solution:
Let ABC be a triangle.
If possible, let A = 90° and B = 90°.
We know that, the sum of the three internal angles of a triangle is 180°.
i.e., A + B + C = 180°
90° + 90° +C = 180°
C = 0° which is a contradiction.
Hence, there cannot be two right angles in a triangle.
Applicafication to switching networks:
Now we consider the statements p and p1 as switches with the property that if one is on, then the other is
off and vice-versa.
Further, a switch allows only two possibilities. They are
(i) it is either open (F) in which case there is no flow of current.
(or)
(ii) it is closed (T) in which case there is a flow of current. Hence, every switch has only two truth values
T or F
Let p and q denote two switches. We can connect p and q by using a wire in a series or parallel combina-
tion as shown below.
p q
B
A
series combination
p
A B
q
parallel combination
Note:
Let p q denote the series combination and p q denote the parallel combination.
Switching network:
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MATHEMATICS
A switching network is an arrangement of wires and switches in series and parallel combinations.
So, such a network can be described by using the connectives and .
Example:
1. Describe the behavior of flow of current from A to B in the following circuit network.
Solution:
p q
A B
p׀
The given network can be described by the compound statement (p q) p. Truth table of (pq) p
is:
p q p׀ pq (p q) p׀
T T F T T
T F F F F
F T T F T
F F T F T
So, current flows from A to B if
(i) p is closed, q is closed,
(ii) p is open, q is closed and
(iii) p is open, q is open.
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MATHEMATICAL
10 INDUCTION
456
MATHEMATICS
Introduction
The process of mathematical induction is a indirect method which helps us to prove complex mathematical
formulae, that cannot be easily proved by direct methods.
For example, to prove that ‘n(n + 1) is always divisible by 2’ for n being a natural number, we can
substitute n = 1, 2, 3…… in n(n + 1), and check in each case if the result is divisible by 2. After checking, for a
few of values, we can say that the formula is likely to be correct. Since, we cannot substitute all possible values
of n, to prove the formula we use the principle of mathematical induction to prove the given formula.
The principle of mathematical induction:
If P(n) is a statement such that,
(i) P(n) is true for n =1
(ii) P(n) is true for n = k + 1, when it is true for n = k, where k is a natural number
then the statement P(n) is true for all natural numbers.
Let us prove some results using this principle
Example 1:
n( n 1)
Prove that 1 + 2 + 3 + …. + n = .
2
Solution:
n( n 1)
Let P(n): 1+2+….+n = be the given statement.
2
Step 1: Put n = 1
1(1 1)
Then, L.H.S. = 1 and R.H.S. = = 1.
2
L.H.S. = R.H.S.
P(n) is true for n = 1.
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Step 2: Assume that P(n) is true for n = k.
k k 1
1 + 2 + 3 + …. + k =
2
Adding (k+1) on both sides, we get
k( k 1)
1 + 2 + 3 + …. + k + (k + 1) = + (k + 1)
2
k
= (k + 1) 1
2
( k 1)( k 2)
=
2
( k 1)( k 1 1)
=
2
P(n) is true for n = k + 1
By the principle of mathematical induction P(n) is true for all natural numbers n.
n( n 1)
Hence, 1 + 2 + 3 + …. + n = for all nN
2
Example 2:
Prove that 1 + 3 + 5 + …. + (2n – 1) = n2
Solution:
Let P(n): 1 + 3 + 5 + …. + (2n – 1) = n2 be the given statement
Step 1: Put n =1
Then, L.H.S. = 1
R.H.S. = (1)2 = 1
L.H.S. = R.H.S.
P(n) is true for n = 1.
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MATHEMATICS
Step 2: Assume that P(n) is true for n = k.
1 + 3 + 5 + …. + (2k – 1) = k2
Adding 2k+1 on both sides, we get
1 + 3 + 5 +….+ (2k – 1) + (2k + 1) = k2 + (2k + 1) = (k + 1)2
1 + 3 + 5 + …. + (2k – 1) + (2(k + 1) – 1) = (k + 1)2
P(n) is true for n = k + 1.
By the principle of mathematical induction P(n) is true for all natural numbers ‘n’.
Hence, 1 + 3 + 5 + …. + (2n – 1) = n2, for all nN
Example 3:
n ( n 1)( n 2)
Prove that 1.2 + 2.3 + 3.4 +…+ n.(n + 1) =
3
Solution:
n( n 1)(n 2)
Let P(n): 1.2 + 2.3 + 3.4 +…+ n.(n + 1) = be the given statement
3
Step 1: Put n = 1
Then, L.H.S.= 1.2 = 2
1(1 1)(1 2) 23
R.H.S. = = =2
3 3
L.H.S. = R.H.S.
P(n) is true for n = 1.
Step 2: Assume that P(n) is true for n = k.
k ( k 1)( k 2)
1.2 + 2.3 + 3.4 + …. + k(k + 1) =
3
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Adding (k + 1)(k + 2) on both sides, we get 1.2 + 2.3 + 3.4 + ……. +
k(k + 1) + (k + 1)(k + 2)
k( k 1)( k 2)
= ( k 1)( k 2)
3
k
= (k+1)(k+2) 1
3
( k 1)( k 2)( k 3)
=
3
1.2 + 2.3 + 3.4 + …. + k.(k + 1) + (k + 1)(k + 2)
( k 1)( k 1 1)(k 1 2)
=
3
P(n) is true for n = k + 1.
By the principle of mathematical induction P(n) is true for all natural numbers
n ( n 1)( n 2)
Hence, 1.2 + 2.3 + 3.4 + … + n.(n + 1) = , nN
3
Example 4:
Prove that 3n+1 > 3(n + 1)
Solution:
Let P(n) : 3n+1> 3(n + 1)
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MATHEMATICS
Step 1: Put n = 1
Then, 32 > 3(2)
p(n) is true for n = 1
Step 2: Assume that P(n) is true for n = k
Then, 3k+1 > 3(k + 1)
Multiplying throughout with ‘3’.
3k+1× 3 > 3(k + 1) × 3 = 9k + 9 = 3( k + 2) + (6k + 3) > 3(k + 2)
3k 11 3( k 1 1)
P(n) is true for n = k + 1
By the principle of mathematical induction, P(n) is true for all n Î N.
Hence, 3n+1 > 3(n + 1), “n N
Example 5:
Prove that 7 is a factor of 23n – 1 for all natural numbers n.
Solution:
Let P(n) : 7 is a factor of 23n – 1 be the given statement
Step 1: When n = 1,
23(1) – 1 = 7 and 7 is a factor of itself.
P(n) is true for n = 1
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Step 2: Let P(n) be true for n = k.
7 is a factor of 23k – 1.
23k – 1 = 7M, where MÎN.
23K = 7M + 1 —————(1)
Now consider 23(k+1) – 1 = 23k+3 – 1 = 23k.23 – 1
= 8(7M+1) – 1 (using (1)) = 56M + 7 (As 23k = 7m + 1)
23(k+1) – 1 = 7(8M + 1)
7 is a factor of 23(k+1) – 1
P(n) is true for n = k + 1
By the principle of mathematical induction, P(n) is true for all natural numbers n.
Hence, 7 is a factor 23n –1 for all nN.
462
MATHEMATICS
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11 BINOMIAL THEOREM
464
MATHEMATICS
Introduction
Binomial expression: An algebraic expression containing only two terms is called a binomial expression.
For example, x + 2y, 3x + 5y, 8x – 7y etc,
We know that, (a + b)2 = a2 + 2ab + b2.
(a + b)3 = (a + b) (a + b)2
= a3 + 3a2b + 3ab2 + b3
Now using a similar approach we can arrive at the expressions for (a + b)4, (a + b)5 etc. However, when
the index is large, this process becomes very cumbersome. Hence, we need a simpler method to arrive at the
expression for (a + b)n,
for n = 1, 2, 3…..
The binomial theorem is the appropriate tool in this case. It helps us arrive at the expression for (a + b)n, for
any value of n, by using a few standard coefficients also know as binomial coefficients.
Now, consider the following cases in which we find the expansions when a binomial expression is raised to
different powers.
(x + y)1 = x + y
(x + y)2 = x2 + 2xy + y2
(x + y)3 = x3 + 3x2y + 3xy2 + y3
(x + y)4 = x4 + 4x3y + 6x2y2 + 4xy3 + y4
In the above examples, the coefficients of the variables in the expansions of the powers of the binomial
expression are called binomial coefficients.
When the binomial coefficients are listed, for different values of n, we see a definite pattern being followed.
This pattern is given by the Pascal Triangle.
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Pascal Triangle
The exponent in the The coefficients of the terms in
binomial the expansion
1 1 1
1 2 1
2
1 3 3 1
3
1 4 6 4 1
4
1 5 10 10 5 1
5
This definite pattern, show above, can be used to write the binomial expansions for higher powers such as
n = 6, 7, 8…. so on. The binomial theorem gives us a general algebraic formula by means of which any power of
a binomial expression can be expanded into a series of simpler terms.
Before we take up the binomial theorem, let us review the concepts of factorial notation and the nCr
representation.
Factorial notation and nCr representation:
The factorial of n is denoted by n! and is defined as n! = 1 2 3 …... (n – 1) n
For example, 4! = 1 2 3 4 and 6! = 1 2 3 4 5 6.
Also, 0! = 1 and n! = n (n – 1)!.
n!
For 0 d” r d” n, we define nCr as nCr = n r ! r!
6! 6 5 4!
for example, 6C2 = 6 2 ! 2! = 4! 2! = 15
466
MATHEMATICS
also, nC0 = nCn = 1; nC1 = nCn–1 = n and nCr = nCn–r
for example, 10C2 = 10C8 and if nC3 = nC5, then n = 3 + 5 = 8
Binomial theorem:
If n is a positive integer,
(x + y)n = nC0 xn + nC1 xn–1y + nC2 xn–2 y2 + ……... + nCr xn–r yr + …….. + nCn yn
Important inferences from the above expansion:
The number of terms in the expansion is n + 1.
The exponent of x goes on decreasing by ‘1’ from left to right and the power of ‘y’ goes on increasing by
1 from left to right.
In each term of the expansion the sum of the exponents of x and y is equal to the exponent (n) of the
binomial expression.
the coefficients of the terms that are equidistant from the beginning and the end have numerically equal, i.e.,
n
C0 = nCn; nC1 = nCn–1; nC2 = nCn–2 and so on.
The general term in the expansion of (x + y)n is given by Tr + 1 = nCr xn–r yr.
On substituting ‘– y’ in place of ‘y’ in the expansion, we get
(x–y)n = nC0 xn – nC1 xn–1 y + nC2 xn–2 y2 – nC3 xn–3 y3 + … + (–1)n nCnyn
The general term in the expansion (x – y)n is Tr + 1 = (–1)r nCr xn–r yr.
Example 1:
Expand (x + 2y)5.
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Solution:
(x + 2y)5 = 5C0 x5 + 5C1 x5–1 (2y) + 5C2 x5–2 (2y)2 + 5C3 x5–3 (2y)3
+ 5C4 x5–4 (2y)4 + 5C5 (2y)5.
(x + 2y)5 = x5 + 5x4 (2y) + 10x3 4y2 + 10x28y3 + 5x 16y4 + 25 y5
= x5 + 10x4y + 40x3y2 + 80x2y3 + 80xy4 + 32y5
Example 2:
Find the 3rd term in the expansion of (3x–5y)7.
Solution:
The general term in (x – y)n is Tr + 1 = (–1)r nCr xn–ryr
T3 = T2 + 1 = (–1)2 7C2 (3x)7–2 (5y)2 = 7C2 (3x)5 (5y)2.
Middle terms in the expansion of (x + y)n:
Depending on the nature of n, i.e., whether n is even or odd, there may exist one or two middle terms.
Case 1:
When n is an even number, then there is only one middle term in the expansion
th
n
(x + y) , which is 1 term.
n
2
Case 2:
th
n 1
n
When n is odd number, there will be two middle terms in the expansion of (x + y) , which are and
2
th
n 3
terms.
2
468
MATHEMATICS
Example 1:
Find the middle term in the expansion of (2x + 3y)8.
Solution:
th
8
Since n is even number, 1 term i.e., 5th term is the middle term in
2
(2x + 3y)8.
T5 = T4 + 1 = 8C4 (2x)8–4 (3y)4 = 8C4 (2x)4 (3y)4
Example 2:
Find the middle terms in the expansion of (5x–7y)7.
Solution:
Since n is an odd number, the expansion contains two middle terms.
th th
7 1 7 3
and terms are the two middle terms in the expansion of (5x –7y)7.
2 2
T4 = T3 + 1 = (–1)3 7C3 (5x)7–3 (7y)3
= –7C3 (5x)4 (7y)3
T5 = T4 + 1 = (–1)4× 7C4 (5x)7–4 × (7y)4
= –7C4 (5x)3 (7y)4
Term independent of x:
n
p 1
In an expansion of form x q , the term for which the exponent of x is 0 is said to be the term that is
x
independent of x or a constant term.
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2
1 1
For example, in the expansion x = x2 + 2 + 2 , the 2nd term is independent of ‘x’.
x x
Example 1:
4
1
Find the term independent of x in x .
x
Solution:
Let Tr + 1 be the term independent term of x in the given expansion.
r
4 4–r
1 x 4 r 4
Tr + 1 = Cr x 4
= Cr = Cr x4–2r
x xr
For the term independent of x the power of x should be zero.
4 – 2r = 0 or r = 2.
T2+1 = T3 term, is the independent term of the expansion.
Note:
If r is not a positive integer, then the expansion does not contain constant term.
Example 2:
6
2 1
Find the coefficient of x in x 3 .
2
x
Solution:
Let Tr + 1 be the term containing x2.
r
1
6
Tr + 1 = Cr (x )2 6–r 3
x
1
= 6Cr x12–2r = 6Cr x12–5r
x 3r
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MATHEMATICS
As the coefficient of x is 2
12 – 5r = 2 r = 2.
Coefficient of x2 = 6C2 = 15.
The greatest coefficient in the expansion of (1 + x)n:
(where n is a positive integer)
The coefficient of the (r + 1)th term in the expansion of (1 + x)n is nCr.
n
Cr is maximum when r = n/2 (if n is even) and
n 1 n 1
r= or (if n is odd)
2 2
Example 1:
Find the total number of terms in the expansion of 2 3x 15 2 3x 15 .
Solution:
(2 + 3x)15
= 15C0(2)15 + 15C1(2)14 (3x)1 + …. + 15C14(2)1 (3x)14 + 15C15(3x)15 and (2 3x)15
= 15C0(2)15 15C1(2)14 (3x1) + ……. + 15C14(2)1 (3x)14 15C15(3x)15
Adding the two equations, we see that the terms in even positions get cancelled, and we get
(2 + 3x)15 + (2 3x)15 = 2[15C0(2)15 + 15C2(2)13 (3x)2 + …. + 15C14(2)1 (3x)14]
Total number of terms = 8.
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n 1
Alternately, the number of terms in (a + x)n + (a – x)n, if n is odd is . Hence in this case, the number
2
15 1
of terms are = 8.
2
Example 2:
n
2 1
If the expansion x 3 is to contain an independent term, then what should be the value of n?
x
Solution:
General term, Tr+1 = n Cr x n r y r , for (x + y)n
n
2 1 1
general term of x is
n
C r x 2 n 2 r 3r = n C x 2 n 5 r
x3 x r
For a term to be independent of x, 2n - 5r should be equal to zero,
i.e., 2n - 5r = 0.
2
r= n, since r can take only integral values, n has to be a multiple of 5.
5
Example 3:
9 9
1 1
If the coefficient of x in ax and x-7 in bx are equal, find the relation between a and b?
7
x x
472
MATHEMATICS
Solution:
9 r
1 1
For ax , Tr+1 = 9Cr(ax)9 - r
x x
= 9Cr(a)9 – r (x)9 – 2r as 9 - 2r = 7, r = 1
Coefficient of x7 is 9C1(a)9 - 1 = 9(a)8
9 r
1 1
Now, for bx , Tr+1 = 9Cr(bx)9 – r x
x
= 9Cr(b)9 - r (-1)r x9 - 2r as 9 - 2r = -7, r = 8.
Coefficient of x-7 is 9C8 b9 - 8 (-1)8 = 9b
9a8 = 9b i.e., a8 - b = 0
Example 4:
Find the term independent of ‘x’ in the expansion of
4
1
(1 + x2)4 1 2 .
x
Solution:
4
2 4 1
(1 x ) 1 2
x
= (4 C0 4C1x 2 ..... 4 C4 x 8 ) (4 C0 4C1x 2 ..... 4 C 4 x 8 )
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The term independent of ‘x’ is the term containing the coefficient (4C0 4 C0 4 C1 4 C1 ....4 C4 4 C4 ) =
(4C0)2 (4C1)2 (4C2)2 (4C3)2 (4C4)2
= 12 + 42 + 62 +42 + 12 = 70
Example 5:
Find the sum of the co-efficients of the terms of the expansion
(1 + x + 2x2)6.
Solution:
Substituting x = 1, we have (1 + 1 + 2)6, which gives us the sum of the
co-efficients of the terms of the expansion.
Sum = 46
Example 6:
6
1 2 x
Find the value of x, if the fourth term in the expansion of 2 x .2 is 160.
x
Solution:
3
1
6 2 3 x 3
4 term T3 + 1 = C3 2 ( x ) (2 )
th
x
6 C3 (2 x )3 160
i.e., 20 23x = 160
23x = 8 Þ 23x = 23
x=1
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MATHEMATICS
475
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