MATHEMATICAL TRIPOS Part II 2004

List of Courses
Geometry of Surfaces
Graph Theory
Number Theory
Coding and Cryptography
Algorithms and Networks
Computational Statistics and Statistical Modelling
Quantum Physics
Statistical Physics and Cosmology
Symmetries and Groups in Physics
Transport Processes
Theoretical Geophysics
Mathematical Methods
Nonlinear Waves
Markov Chains
Principles of Dynamics
Functional Analysis
Groups, Rings and Fields
Electromagnetism
Dynamics of Differential Equations
Logic, Computation and Set Theory
Principles of Statistics
Stochastic Financial Models
Foundations of Quantum Mechanics
General Relativity
Numerical Analysis
Nonlinear Dynamical Systems
Combinatorics
Representation Theory
Galois Theory
Differentiable Manifolds
Algebraic Topology
Number Fields
Hilbert Spaces
Riemann Surfaces
Algebraic Curves
Probability and Measure
Applied Probability
Information Theory
Optimization and Control
Partial Differential Equations
Methods of Mathematical Physics
Electrodynamics
Statistical Physics
Applications of Quantum Mechanics
Fluid Dynamics II
Waves in Fluid and Solid Media

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A2/7 Geometry of Surfaces

(i) What is a geodesic on a surface M with Riemannian metric, and what are geodesic
polar co-ordinates centred at a point P on M ? State, without proof, formulae for the
Riemannian metric and the Gaussian curvature in terms of geodesic polar co-ordinates.
(ii) Show that a surface with constant Gaussian curvature 0 is locally isometric to the
Euclidean plane.

A3/7 Geometry of Surfaces

(i) The catenoid is the surface C in Euclidean R3 , with co-ordinates x, y, z and
Riemannian metric ds2 = dx2 + dy 2 + dz 2 obtained by rotating the curve y = cosh x
about the x-axis, while the helicoid is the surface H swept out by a line which lies along
the x-axis at time t = 0, and at time t = t0 is perpendicular to the z-axis, passes through
the point (0, 0, t0 ) and makes an angle t0 with the x-axis.
Find co-ordinates on each of C and H and write x, y, z in terms of these co-ordinates.
(ii) Compute the induced Riemannian metrics on C and H in terms of suitable co-
ordinates. Show that H and C are locally isometric. By considering the x-axis in H,
show that this local isometry cannot be extended to a rigid motion of any open subset of
Euclidean R3 .

A4/7 Geometry of Surfaces

Write an essay on the Gauss–Bonnet theorem and its proof.

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A1/8 Graph Theory

(i) Let G be a connected graph of order n ≥ 3 such that for any two vertices x and y,

d(x) + d(y) ≥ k.

Show that if k < n then G has a path of length k, and if k = n then G is Hamiltonian.
(ii) State and prove Hall’s theorem.
[If you use any form of Menger’s theorem, you must state it clearly.]
Let G be a graph with directed edges. For S ⊂ V (G), let

Γ+ (S) = {y ∈ V (G) : xy ∈ E(G) for some x ∈ S} .

Find a necessary and sufficient condition, in terms of the sizes of the sets Γ+ (S), for the
existence of a set F ⊂ E(G) such that at every vertex there is exactly one incoming edge
and exactly one outgoing edge belonging to F .

A2/8 Graph Theory

(i) State a result of Euler, relating the number of vertices, edges and faces of a plane
graph. Show that if G is a plane graph then χ(G) ≤ 5.
(ii) Define the chromatic polynomial pG (t) of a graph G. Show that

pG (t) = tn − a1 tn−1 + a2 tn−2 + . . . + (−1)n an

where a1 , . . . , an are non-negative integers. Explain, with proof, how the chromatic
polynomial is related to the number of vertices, edges and triangles in G. Show that
if Cn is a cycle of length n ≥ 3, then

pCn (t) = (t − 1)n + (−1)n (t − 1).

A4/9 Graph Theory

Write an essay on trees. You should include a proof of Cayley’s result on the
number of labelled trees of order n.
Let G be a graph of order n ≥ 2. Which of the following statements are equivalent
to the statement that G is a tree? Give a proof or counterexample in each case.
(a) G is acyclic and e(G) ≥ n − 1.
(b) G is connected and e(G) ≤ n − 1.
(c) G is connected, triangle-free and has at least two leaves.
(d) G has the same degree sequence as T , for some tree T .

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A1/9 Number Theory

(i) State the law of quadratic reciprocity. For p 6= 5 an odd prime, evaluate the
Legendre symbol  
5
.
p

(ii) (a) Let p1 , . . . , pm and q1 , . . . , qn be distinct odd primes. Show that there exists an
integer x that is a quadratic residue modulo each of p1 , . . . , pm and a quadratic non-residue
modulo each of q1 , . . . , qn .
(b) Let p be an odd prime. Show that

p−1  
X a
= 0.
a=1
p

(c) Let p be an odd prime. Using (b) or otherwise, evaluate

p−2    
X a a+1
.
a=1
p p
   
x2 y y
[Hint for (c): Use the equality p = p , valid when p does not divide x.]

A3/9 Number Theory

(i) Find a solution in integers of the Pell equation x2 − 17y 2 = 1.
(ii) Define the continued fraction expansion of a real number θ > 1 and show that it
converges to θ.
Show that if N > 0 is a nonsquare√integer and x and y are integer solutions of
x − N y 2 = 1, then x/y is a convergent of N .
2

A4/10 Number Theory

Write an essay on pseudoprimes and their role in primality testing. You should
discuss pseudoprimes, Carmichael numbers, and Euler and strong pseudoprimes. Where
appropriate, your essay should include small examples to illustrate your statements.

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A1/10 Coding and Cryptography

(i) What is a linear code? What does it mean to say that a linear code has length n
and minimum weight d? When is a linear code perfect? Show that, if n = 2r − 1, there
exists a perfect linear code of length n and minimum weight 3.
(ii) Describe the construction of a Reed-Muller code. Establish its information rate
and minimum weight.

A2/9 Coding and Cryptography

(i) Describe how a stream cypher operates. What is a one-time pad?
A one-time pad is used to send the message x1 x2 x3 x4 x5 x6 y7 which is encoded as
0101011. By mistake, it is reused to send the message y0 x1 x2 x3 x4 x5 x6 which is encoded
as 0100010. Show that x1 x2 x3 x4 x5 x6 is one of two possible messages, and find the two
possibilities.
(ii) Describe the RSA system associated with a public key e, a private key d and the
product N of two large primes.
Give a simple example of how the system is vulnerable to a homomorphism attack.
Explain how a signature system prevents such an attack. [You are not asked to give an
explicit signature system.]
Explain how to factorise N when e, d and N are known.

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A2/10 Algorithms and Networks

(i) Define the minimum path and the maximum tension problems for a network with
span intervals specified for each arc. State without proof the connection between the two
problems, and describe the Max Tension Min Path algorithm of solving them.
(ii) Find the minimum path between nodes S and S0 in the network below. The span
intervals are displayed alongside the arcs.

[−5, 5] [−1, 2]

[0, 2]
1 2

[−6, 2] [−3, 3]

[−2, 3] 3 [0, 4]

[0, ∞) [−∞, 2]

[−6, 0]
4 5

[−1, 11]
(−∞, ∞)

[−1, 3] S0 [−2, 0]

[−1, 9]
[−1, 6]

6 7

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A3/10 Algorithms and Networks

(i) Consider the problem
minimise f (x)
(∗)
subject to g(x) = b, x ∈ X,
where f : Rn −→ R, g : Rn −→ Rm , X ⊆ Rn , x ∈ Rn and b ∈ Rm . State the Lagrange
Sufficiency Theorem for problem (∗). What is meant by saying that this problem is strong
Lagrangian? How is this related to the Lagrange Sufficiency Theorem? Define a supporting
hyperplane and state a condition guaranteeing that problem (∗) is strong Lagrangian.
(ii) Define the terms flow, divergence, circulation, potential and differential for a
network with nodes N and arcs A.
State the feasible differential problem for a network with span intervals D(j) =
[d− (j), d+ (j)], j ∈ A.
State, without proof, the Feasible Differential Theorem.
[You must carefully define all quantities used in your statements.]
Show that the network below does not support a feasible differential.

[1, 4]

[−1, 1] [2, 6]

[−2, 2] [1, 5]

[0, 2] [0, 6] [1, 3]

[0, 6]

[−1, 1] [−5, 2]

[−1, 6] [−6, 1] [2, 7]

[−1, 7] [3, 9]

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A4/11 Algorithms and Networks

(i) Consider an unrestricted geometric programming problem

min g(t), t = (t1 , . . . , tm ) > 0 , (∗)

where g(t) is given by

n
X
g(t) = ci t1ai1 . . . tamim
i=1

with n ≥ m and positive coefficients c1 . . . , cn . State the dual problem of (∗) and show
that if λ∗ = (λ∗1 , . . . , λ∗n ) is a dual optimum then any positive solution to the system

1 ∗
ta1 i1 . . . tamim = λ v(λ∗ ), i = 1, . . . , n ,
ci i

gives an optimum for primal problem (∗). Here v(λ) is the dual objective function.
(ii) An amount of ore has to be moved from a pit in an open rectangular skip which
is to be ordered from a supplier.
The skip cost is £36 per 1m2 for the bottom and two side walls and £18 per 1m2
for the front and the back walls. The cost of loading ore into the skip is £3 per 1m3 , the
cost of lifting is £2 per 1m3 , and the cost of unloading is £1 per 1m3 . The cost of moving
an empty skip is negligible.
Write down an unconstrained geometric programming problem for the optimal size
(length, width, height) of skip minimizing the cost of moving 48m3 of ore. By considering
the dual problem, or otherwise, find the optimal cost and the optimal size of the skip.

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A1/13 Computational Statistics and Statistical Modelling

(i) Assume that the n-dimensional vector Y may be written as Y = Xβ + , where X
is a given n × p matrix of rank p, β is an unknown vector, and

 ∼ Nn (0, σ 2 I).

Let Q(β) = (Y − Xβ)T (Y − Xβ). Find β̂, the least-squares estimator of β, and state
without proof the joint distribution of β̂ and Q(β̂).
(ii) Now suppose that we have observations (Yij , 1 6 i 6 I, 1 6 j 6 J) and consider
the model
Ω : Yij = µ + αi + βj + ij ,
where (αi ), (βj ) are fixed parameters with Σαi = 0, Σβj = 0, and (ij ) may be assumed
independent normal variables, with ij ∼ N (0, σ 2 ), where σ 2 is unknown.
(a) Find (α̂i ), (β̂j ), the least-squares estimators of (αi ), (βj ).
(b) Find the least-squares estimators of (αi ) under the hypothesis H0 : βj = 0 for
all j.
(c) Quoting any general theorems required, explain carefully how to test H0 ,
assuming Ω is true.
(d) What would be the effect of fitting the model Ω1 : Yij = µ + αi + βj + γij + ij ,
where now (αi ), (βj ), (γij ) are all fixed unknown parameters, and (ij ) has the distribution
given above?

A2/12 Computational Statistics and Statistical Modelling

(i) Suppose we have independent observations Y1 , . . . , Yn , and we assume that for
i = 1, . . . , n, Yi is Poisson with mean µi , and log(µi ) = β T xi , where x1 , . . . , xn are given
covariate vectors each of dimension p, where β is an unknown vector of dimension p,
and p < n. Assuming that {x1 , . . . , xn } span Rp , find the equation for β̂, the maximum
likelihood estimator of β, and write down the large-sample distribution of β̂.
(ii) A long-term agricultural experiment had 90 grassland plots, each 25m × 25m,
differing in biomass, soil pH, and species richness (the count of species in the whole plot).
While it was well-known that species richness declines with increasing biomass, it was not
known how this relationship depends on soil pH, which for the given study has possible
values “low”, “medium” or “high”, each taken 30 times. Explain the commands input, and
interpret the resulting output in the (slightly edited) R output below, in which “species”
represents the species count.
(The first and last 2 lines of the data are reproduced here as an aid. You may
assume that the factor pH has been correctly set up.)

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> species
pH Biomass Species
1 high 0.46929722 30
2 high 1.73087043 39
.......................
.......................
89 low 4.36454121 7
90 low 4.87050789 3

> summary(glm(Species ~Biomass, family = poisson))

Call:
glm(formula = Species ~ Biomass, family = poisson)

Coefficients:
Estimate Std. Error z value Pr(>|z|)
(Intercept) 3.184094 0.039159 81.31 < 2e-16
Biomass -0.064441 0.009838 -6.55 5.74e-11

(Dispersion parameter for poisson family taken to be 1)

Null deviance: 452.35 on 89 degrees of freedom

Residual deviance: 407.67 on 88 degrees of freedom

Number of Fisher Scoring iterations: 4

> summary(glm(Species ~pH*Biomass, family = poisson))

Call:
glm(formula = Species ~ pH * Biomass, family = poisson)

Coefficients:
Estimate Std. Error z value Pr(>|z|)
(Intercept) 3.76812 0.06153 61.240 < 2e-16
pHlow -0.81557 0.10284 -7.931 2.18e-15

Question continues on next page.

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pHmid -0.33146 0.09217 -3.596 0.000323

Biomass -0.10713 0.01249 -8.577 < 2e-16
pHlow:Biomass -0.15503 0.04003 -3.873 0.000108
pHmid:Biomass -0.03189 0.02308 -1.382 0.166954

(Dispersion parameter for poisson family taken to be 1)

Null deviance: 452.346 on 89 degrees of freedom

Residual deviance: 83.201 on 84 degrees of freedom

Number of Fisher Scoring iterations: 4

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A4/14 Computational Statistics and Statistical Modelling

Suppose that Y1 , . . . , Yn are independent observations, with Yi having probability
density function of the following form
 
yi θi − b(θi )
f (yi |θi , φ) = exp + c(yi , φ)
φ

where E(Yi ) = µi and g(µi ) = β T xi . You should assume that g( ) is a known function,
and β, φ are unknown parameters, with φ > 0, and also x1 , . . . , xn are given linearly
independent covariate vectors. Show that

∂` X (yi − βi )
= xi ,
∂β g 0 (µi )Vi

where ` is the log-likelihood and Vi = var (Yi ) = φb00 (θi ).

Discuss carefully the (slightly edited) R output given below, and briefly suggest
another possible method of analysis using the function glm ( ).

> s <- scan()

1: 33 63 157 38 108 159
7:
Read 6 items
> r <- scan()
1: 3271 7256 5065 2486 8877 3520
7:
Read 6 items
> gender <- scan(,"")
1: b b b g g g
7:
Read 6 items
> age <- scan(,"")
1: 13&under 14-18 19&over
4: 13&under 14-18 19&over
7:
Read 6 items
> gender <- factor(gender) ; age <- factor(age)
> summary(glm(s/r ~ gender + age,binomial, weights=r))

Coefficients:
Question continues on next page.

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Estimate Std.Error z-value Pr(>|z|)

(Intercept) -4.56479 0.12783 -35.710 < 2e-16
genderg 0.38028 0.08689 4.377 1.21e-05
age14-18 -0.19797 0.14241 -1.390 0.164
age19&over 1.12790 0.13252 8.511 < 2e-16

(Dispersion parameter for binomial family taken to be 1)

Null deviance: 221.797542 on 5 degrees of freedom

Residual deviance: 0.098749 on 2 degrees of freedom

Number of Fisher Scoring iterations: 3

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A1/14 Quantum Physics

(i) Each particle in a system of N identical fermions has a set of energy levels Ei with
degeneracy gi , where i = 1, 2, . . .. Derive the expression
gi
N̄i = ,
eβ(Ei −µ) + 1

for the mean number of particles N̄i with energy Ei . Explain the physical significance of
the parameters β and µ.
(ii) The spatial eigenfunctions of energy for an electron of mass m moving in two
dimensions and confined to a square box of side L are

2  n πx 
1
 n πy 
2
ψn1 n2 (x) = sin sin ,
L L L

where ni = 1, 2, . . . (i = 1, 2). Calculate the associated energies.

Hence show that when L is large the number of states in energy range E → E + dE is

mL2
dE .
2π~2
How is this formula modified when electron spin is taken into account?
The box is filled with N electrons in equilibrium at temperature T . Show that the
chemical potential µ is given by

1  2

µ = log eβπ~ ρ/m − 1 ,
β

where ρ is the number of particles per unit area in the box.

What is the value of µ in the limit T → 0?
Calculate the total energy of the lowest state of the system of particles as a function of N
and L.

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A2/14 Quantum Physics

(i) A simple model of a crystal consists of an infinite linear array of sites equally
spaced with separation b. The probability amplitude for an electron to be at the n-th site
is cn , n = 0, ±1, ±2, . . .. The Schrödinger equation for the {cn } is

Ecn = E0 cn − A(cn−1 + cn+1 ) ,

where A is real and positive. Show that the allowed energies E of the electron must lie in
a band |E − E0 | ≤ 2A, and that the dispersion relation for E written in terms of a certain
parameter k is given by
E = E0 − 2A cos kb .

What is the physical interpretation of E0 , A and k?

(ii) Explain briefly the idea of group velocity and show that it is given by

1 dE(k)
v = ,
~ dk

for an electron of momentum ~k and energy E(k).

An electron of charge q confined to one dimension moves in a periodic potential
under the influence of an electric field E. Show that the equation of motion for the electron
is
qE d2 E
v̇ = 2 ,
~ dk 2
where v(t) is the group velocity of the electron at time t. Explain why
−1
d2 E

∗ 2
m = ~
dk 2

can be interpreted as an effective mass.

Show briefly how the absence from a band of an electron of charge q and effective
mass m∗ < 0 can be interpreted as the presence of a ‘hole’ carrier of charge −q and
effective mass −m∗ .
In the model of Part (i) show that
(a) for k 2  12/b2 an electron behaves like a free particle of mass ~2 /(2Ab2 );
(b) for (π/b − k)2  12/b2 a hole behaves like a free particle of mass ~2 /(2Ab2 ).

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A4/16 Quantum Physics

Explain the operation of the np junction. Your account should include a discussion
of the following topics:
(a) the rôle of doping and the fermi-energy;
(b) the rôle of majority and minority carriers;
(c) the contact potential;
(d) the relationship I(V ) between the current I flowing through the junction and the
external voltage V applied across the junction;
(e) the property of rectification.

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A1/16 Statistical Physics and Cosmology

(i) Consider a homogeneous and isotropic universe with mass density ρ(t), pressure
P (t) and scale factor a(t). As the universe expands its energy E decreases according
to the thermodynamic relation dE = −P dV where V is the volume. Deduce the fluid
conservation law
ȧ  P
ρ̇ = −3 ρ + 2 .
a c
Apply the conservation of total energy (kinetic plus gravitational potential) to a test
particle on the edge of a spherical region in this universe to obtain the Friedmann equation
 2
ȧ 8πG kc2
= ρ− 2 ,
a 3 a

where k is a constant. State clearly any assumptions you have made.

(ii) Our universe is believed to be flat (k = 0) and filled with two major components:
pressure-free matter (PM = 0) and dark energy with equation of state PQ = −ρQ c2 where
the mass densities today (t = t0 ) are given respectively by ρM0 and ρQ0 . Assume that
each component independently satisfies the fluid conservation equation to show that the
total mass density can be expressed as
ρM0
ρ(t) = + ρQ0 ,
a3

where we have set a(t0 ) = 1.

Now consider the substitution b = a3/2 in the Friedmann equation to show that the
solution for the scale factor can be written in the form

a(t) = α(sinh βt)2/3 ,

where α and β are constants. Setting a(t0 ) = 1, specify α and β in terms of ρM0 , ρQ0 and
t0 . Show that the scale factor a(t) has the expected behaviour for an Einstein-de Sitter
universe at early times (t → 0) and that the universe accelerates at late times (t → ∞).
√
[Hint: Recall that dx/ x2 + 1 = sinh−1 x .]
R

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A3/14 Statistical Physics and Cosmology

(i) In equilibrium, the number density of a non-relativistic particle species is given by
 3/2
2πmkT 2
n = gs e(µ−mc )/kT
,
h2

where m is the mass, µ is the chemical potential and gs is the spin degeneracy. At around
t = 100 seconds, deuterium D forms through the nuclear fusion of nonrelativistic protons
p and neutrons n via the interaction:

p+n ↔ D.

What is the relationship between the chemical potentials of the three species when they are
in chemical equilibrium? Show that the ratio of their number densities can be expressed
as 3/2
h2

nD
≈ eBD /kT ,
nn np πmp kT
where the deuterium binding energy is BD = (mn +mp −mD )c2 and you may take gD = 4.
Now consider the fractional densities Xa = na /nB , where nB is the baryon number of the
universe, to re-express the ratio above in the form

XD
Xn Xp

which incorporates the baryon-to-photon ratio η of the universe. [You may assume that
the photon density is nγ = 16πζ(3) 3
(hc)3 (kT ) .] From this expression, explain why deuterium
does not form until well below the temperature kT ≈ BD .
(ii) The number density n = N/V for a photon gas in equilibrium is given by the
formula
8π ∞ ν 2 dν
Z
n= 3 ,
c 0 ehν/kT − 1
where ν is the photon frequency. By considering the substitution x = hν/kT , show that
the photon number density can be expressed in the form

n = αT 3 ,

where the constant α need not be evaluated explicitly.

State the equation of state for a photon gas and explain why the chemical potential
of the photon vanishes. Assuming that the photon energy density  = E/V = (4σ/c)T 4 ,
use the first law dE = T dS − P dV + µdN to show that the entropy density is given by

16σ 3
s = S/V = T .
3c
Hence explain why, when photons are in equilibrium at early times in our universe, their
temperature varies inversely with the scale factor: T ∝ a−1 .

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A4/18 Statistical Physics and Cosmology

(a) Consider an ideal gas of Fermi particles obeying the Pauli exclusion principle
with a set of one-particle energy eigenstates Ei . Given the probability pi (ni ) at tempera-
ture T that there are ni particles in the eigenstate Ei :

e(µ−Ei )ni /kT

pi (ni ) = ,
Zi
determine the appropriate normalization factor Zi . Use this to find the average number
n̄i of Fermi particles in the eigenstate Ei .
Explain briefly why in generalizing these discrete eigenstates to a continuum in
momentum space (in the range p to p + dp) we must multiply by the density of states
4πgs V 2
g(p)dp = p dp ,
h3
where gs is the degeneracy of the eigenstates and V is the volume.
(b) With the energy expressed as a momentum integral
Z ∞
E= E(p)n̄(p)dp ,
0

consider the effect of changing the volume V so slowly that the occupation numbers do not
change (i.e. particle number N and entropy S remain fixed). Show that the momentum
varies as dp/dV = −p/3V and so deduce from the first law expression
 
∂E
= −P
∂V N,S

that the pressure is given by

Z ∞
1
P = pE 0 (p)n̄(p)dp .
3V 0

Show that in the non-relativistic limit P = 32 U/V where U is the internal energy, while
for ultrarelativistic particles P = 31 E/V .
(c) Now consider a Fermi gas in the limit T → 0 with all momentum eigenstates
filled up to the Fermi momentum pF . Explain why the number density can be written as
4πgs pF 2
Z
n= 3 p dp ∝ p3F .
h 0

From similar expressions for the energy, deduce in both the non-relativistic and ultra-
relativistic limits that the pressure may be written as

P ∝ nγ ,

where γ should be specified in each case.

(d) Examine the stability of an object of radius R consisting of such a Fermi
degenerate gas by comparing the gravitational binding energy with the total kinetic energy.
Briefly point out the relevance of these results to white dwarfs and neutron stars.

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A1/19 Symmetries and Groups in Physics

(i) State and prove Maschke’s theorem for finite-dimensional representations of finite
groups.
(ii) S3 is the group of bijections on {1, 2, 3}. Find the irreducible representations of
S3 , state their dimensions and give their character table.
Let T2 be the set of objects T2 = {ai1 i2 : i1 , i2 = 1, 2, 3}. The operation of the
permutation group S3 on T2 is defined by the operation of the elements of S3 separately
on each index i1 and i2 . For example,

P12 : a13 → a23 , P231 : a23 → a31 , P13 : a33 → a11 .

By considering a representative operator from each conjugacy class of S3 , find the

table of group characters for the representation T2 of S3 acting on T2 . Hence, deduce the
irreducible representations into which T2 decomposes.

A3/15 Symmetries and Groups in Physics

(i) Show that the character of an SU (2) transformation in the 2l + 1 dimensional
irreducible representation dl is given by

sin [(l + 1/2)θ]

χl (θ) = .
sin [θ/2]

What are the characters of irreducible SO(3) representations?

(ii) The isospin representation of two-particle states of pions and nucleons is spanned
by the basis T = {|π + pi, |π + ni, |π 0 pi, |π 0 ni, |π − pi, |π − ni}.
Pions form an isospin triplet with π + = |1, 1i, π 0 = |1, 0i, π − = |1, −1i; and
nucleons form an isospin doublet with p = |1/2, 1/2i, n = |1/2, −1/2i. Find the values of
the isospin for the irreducible representations into which T will decompose.
p
Using I− |j, mi = (j − m + 1)(j + m) |j, m − 1i, write the states of the basis T in
terms of isospin states.
Consider the transitions
π+ p → π+ p
π− p → π− p
π− p → π0 n
and show that their amplitudes satisfy a linear relation.

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A1/18 Transport Processes

(i) In an experiment, a finite amount M of marker gas of diffusivity D is released at
time t = 0 into an infinite tube in the neighbourhood of the origin x = 0. Starting from
the one-dimensional diffusion equation for the concentration C(x, t) of marker gas,

Ct = DCxx ,

use dimensional analysis to show that

M
C= f (ξ)
(Dt)1/2

for some dimensionless function f of the similarity variable ξ = x/(Dt)1/2 .

Write down the equation and boundary conditions satisfied by f (ξ).
(ii) Consider the experiment of Part (i). Find f (ξ) and sketch your answer in the form
of a plot of C against x at a few different times t.
Calculate C(x, t) for a second experiment in which the concentration of marker gas
at x = 0 is instead raised to the value C0 at t = 0 and maintained at that value thereafter.
Show that the total amount of marker gas released in this case becomes greater than M
after a time  2
π M
t= .
16D C0
Show further that, at much larger times than this, the concentration in the first experiment
still remains greater than that in the second experiment for positions x with |x| >
4C0 Dt/M .
Z ∞
2 2 1 2
[Hint: erfc(z) ≡ √ e−u du ∼ √ e−z as z → ∞. ]
π z πz

Part II 2004
 22

A3/16 Transport Processes

(i) Viscous, incompressible fluid of viscosity µ flows steadily in the x-direction in a
uniform channel 0 < y < h. The plane y = 0 is fixed and the plane y = h has constant
x-velocity U . Neglecting gravity, derive from first principles the equations of motion of
the fluid and show that the x-component of the fluid velocity is u(y) and satisfies

0 = −Px + µuyy , (1)

where P (x) is the pressure in the fluid. Write down the boundary conditions on u. Hence
Rh
show that the volume flow rate Q = 0 u dy is given by

U h Px h 3
Q= − . (2)
2 12µ

(ii) A heavy rectangular body of width L and infinite length (in the z-direction) is
pivoted about one edge at (x, y) = (0, 0) above a fixed rigid horizontal plane y = 0. The
body has weight W per unit length in the z-direction, its centre of mass is distance L/2
from the pivot, and it is falling under gravity towards the fixed plane through a viscous,
incompressible fluid. Let α(t)  1 be the angle between the body and the plane. Explain
the approximations of lubrication theory which permit equations (1) and (2) of Part (i)
to apply to the flow in the gap between the two surfaces.
Deduce that, in the gap,
6µα̇
Px = ,
xα3
where α̇ = dα/dt. By taking moments about (x, y) = (0, 0), deduce that α(t) is given by

1 1 2W t
2
− 2 = ,
α α0 3µL

where α(0) = α0 .

Part II 2004
 23

A4/19 Transport Processes

(a) Solute diffuses and is advected in a moving fluid. Derive the transport equation
and deduce that the solute concentration C(x, t) satisfies the advection–diffusion equation

Ct + ∇ · (uC) = ∇ · (D∇C),

where u is the velocity field and D the diffusivity. Write down the form this equation
takes when ∇ · u = 0, both u and ∇C are unidirectional, in the x-direction, and D is a
constant.
(b) A solution occupies the region x > 0, bounded by a semi-permeable membrane at
x = 0 across which fluid passes (by osmosis) with velocity

u = −k (C1 − C(0, t)) ,

where k is a positive constant, C1 is a fixed uniform solute concentration in the region

x < 0, and C(x, t) is the solute concentration in the fluid. The membrane does not allow
solute to pass across x = 0, and the concentration at x = L is a fixed value CL (where
C1 > CL > 0).
Write down the differential equation and boundary conditions to be satisfied by C
in a steady state. Make the equations non-dimensional by using the substitutions

xkC1 C(x) CL
X= , θ(X) = , θL = ,
D C1 C1

and show that the concentration distribution is given by

θ(X) = θL exp [(1 − θ0 )(Λ − X)] ,

where Λ and θ0 should be defined, and θ0 is given by the transcendental equation

θ0 = θL eΛ−Λθ0 . (∗)

What is the dimensional fluid velocity u, in terms of θ0 ?

(c) Show that if, instead of taking a finite value of L, you had tried to take L infinite,
then you would have been unable to solve for θ unless θL = 0, but in that case there would
be no way of determining θ0 .
(d) Find asymptotic expansions for θ0 from equation (∗) in the following limits:
(i) For θL → 0, Λ fixed, expand θ0 as a power series in θL , and equate coefficients
to show that
θ0 ∼ eΛ θL − Λe2Λ θL
2 3


+ O θL .

(ii) For Λ → ∞, θL fixed, take logarithms, expand θ0 as a power series in 1/Λ,

and show that  
log θL 1
θ0 ∼ 1 + +O .
Λ Λ2

What is the limiting value of θ0 in the limits (i) and (ii)?

Question continues on next page.

Part II 2004
 24

(e) Both the expansions in (d) break down when θL = O(e−Λ ). To investigate the
double limit Λ → ∞, θL → 0, show that (∗) can be written as

λ = φeφ

where φ = Λθ0 and λ is to be determined. Show that φ ∼ λ − λ2 + . . . for λ  1, and

φ ∼ log λ − log log λ + . . . for λ  1.
Briefly discuss the implication of your results for the problem raised in (c) above.

Part II 2004
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A1/17 Theoretical Geophysics

(i) What is the polarisation P and slowness s of the time-harmonic plane elastic wave
u = A exp[i(k · x − ωt)]?
Use the equation of motion for an isotropic homogenous elastic medium,

∂2u
ρ = (λ + 2µ)∇(∇ · u) − µ∇ ∧ (∇ ∧ u),
∂t2
to show that s · s takes one of two values and obtain the corresponding conditions on P.
If s is complex show that Re(s) · Im(s) = 0.
(ii) A homogeneous elastic layer of uniform thickness h, S-wave speed β1 and shear
modulus µ1 has a stress-free surface z = 0 and overlies a lower layer of infinite depth,
S-wave speed β2 (> β1 ) and shear modulus µ2 . Show that the horizontal phase speed c
of trapped Love waves satisfies β1 < c < β2 . Show further that
" 1/2 # 1/2
c2 1 − c2 /β22

µ2
tan −1 kh = (1)
β12 µ1 c2 /β12 − 1

where k is the horizontal wavenumber.

Assuming that (1) can be solved to give c(k), explain how to obtain the propagation
speed of a pulse of Love waves with wavenumber k.

Part II 2004
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A2/16 Theoretical Geophysics

(i) Sketch the rays in a small region near the relevant boundary produced by reflection
and refraction of a P -wave incident (a) from the mantle on the core-mantle boundary, (b)
from the outer core on the inner-core boundary, and (c) from the mantle on the Earth’s
surface. [In each case, the region should be sufficiently small that the boundary appears to
be planar.]
Describe the ray paths denoted by SS, P cP , SKS and P KIKP .
Sketch the travel-time (T − ∆) curves for P and P cP paths from a surface source.
(ii) From the surface of a flat Earth, an explosive source emits P -waves downwards into
a stratified sequence of homogeneous horizontal elastic layers of thicknesses h1 , h2 , h3 , . . .
and P -wave speeds α1 < α2 < α3 < . . .. A line of seismometers on the surface records
the travel times of the various arrivals as a function of the distance x from the source.
Calculate the travel times, Td (x) and Tr (x), of the direct wave and the wave that reflects
exactly once at the bottom of layer 1.
Show that the travel time for the head wave that refracts in layer n is given by

n−1 1/2
α2

x X 2hi
Tn = + 1 − 2i .
αn i=1
αi αn

Sketch the travel-time curves for Tr , Td and T2 on a single diagram and show that T2 is
tangent to Tr .
Explain how the αi and hi can be constructed from the travel times of first arrivals
provided that each head wave is the first arrival for some range of x.

Part II 2004
 27

A4/20 Theoretical Geophysics

In a reference frame rotating about a vertical axis with angular velocity f /2,
the horizontal components of the momentum equation for a shallow layer of inviscid,
incompressible, fluid of uniform density ρ are

Du 1 ∂p
− fv = −
Dt ρ ∂x
Dv 1 ∂p
+ fu = − ,
Dt ρ ∂y

where u and v are independent of the vertical coordinate z, and p is given by hydrostatic
balance. State the nonlinear equations for conservation of mass and of potential vorticity
for such a flow in a layer occupying 0 < z < h(x, y, t). Find the pressure p.
By linearising the equations about a state of rest and uniform thickness H, show
that small disturbances η = h − H, where η  H, to the height of the free surface obey

∂2η ∂2η ∂2η

 
− gH + + f 2 η = f 2 η0 − f Hζ0 ,
∂t2 ∂x2 ∂y 2

where η0 and ζ0 are the values of η and the vorticity ζ at t = 0.

Obtain the dispersion relation for homogeneous solutions of the form η ∝ exp[i(kx−
ωt)] and calculate the group velocity of these Poincaré waves. Comment on the form of
these results when ak  1 and ak  1, where the lengthscale a should be identified.
Explain what is meant by geostrophic balance. Find the long-time geostrophically
balanced solution, η∞ and (u∞ , v∞ ), that results from initial conditions η0 = A sgn(x)
and (u, v) = 0. Explain briefly, without detailed calculation, how the evolution from the
initial conditions to geostrophic balance could be found.

Part II 2004
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A2/17 Mathematical Methods

(i) Consider the integral equation

Z b
φ(x) = −λ K(x, t)φ(t)dt + g(x), (†)
a

for φ in the interval a ≤ x ≤ b, where λ is a real parameter and g(x) is given. Describe
the method of successive approximations for solving (†).
Suppose that
|K(x, t)| ≤ M, ∀x, t ∈ [a, b].
By using the Cauchy-Schwarz inequality, or otherwise, show that the successive-approx-
imation series for φ(x) converges absolutely provided

1
|λ| < .
M (b − a)

(ii) The real function ψ(x) satisfies the differential equation

−ψ 00 (x) + λψ(x) = h(x), 0 < x < 1, (?)

where h(x) is a given smooth function on [0, 1], subject to the boundary conditions

ψ 0 (0) = ψ(0), ψ(1) = 0.

By integrating (?), or otherwise, show that ψ(x) obeys

Z 1 Z 1
1 1
ψ(0) = (1 − t)h(t) dt − λ (1 − t)ψ(t) dt.
2 0 2 0

Hence, or otherwise, deduce that ψ(x) obeys an equation of the form (†), with
(
1
2 (1 − x)(1 + t), 0 ≤ t ≤ x ≤ 1,
K(x, t) = 1
2 (1 + x)(1 − t), 0 ≤ x ≤ t ≤ 1,
Z 1
and g(x) = K(x, t)h(t) dt.
0

Deduce that the series solution for ψ(x) converges provided |λ| < 2 .

Part II 2004
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A3/17 Mathematical Methods

(i) Give a brief description of the method of matched asymptotic expansions, as
applied to a differential equation of the type

y 00 + Ky 0 + f (y) = 0, 0 < x < 1,

where 0 <   1, K is a non-zero constant, f is a suitable smooth function and the

boundary values y(0), y(1) are specified. An outline of Van Dyke’s asymptotic matching
principle should be included.
(ii) Consider the boundary-value problem

y 00 + y 0 − (2x + 1)y = 0, y(0) = 0, y(1) = e2

with 0 <   1. Find the integrating factor for the leading-order outer problem. Hence
obtain the first two terms in the outer expansion.
Rewrite the problem using an appropriate stretched inner variable. Hence obtain
the first two terms of the inner exansion.
Use van Dyke’s matching principle to determine all the constants. Hence show that
y (0) = −1 + 25
0
3 + O().

Part II 2004
 30

A4/21 Mathematical Methods

State Watson’s lemma, describing the asymptotic behaviour of the integral
Z A
I(λ) = e−λt f (t) dt, A > 0,
0

as λ → ∞, given that f (t) has the asymptotic expansion

∞
X
f (t) ∼ an tnβ
n=0

as t → 0+ , where β > 0.
Consider the integral
Z b
J(λ) = eλφ(t) F (t)dt,
a

where λ  1 and φ(t) has a unique maximum in the interval [a, b] at c, with a < c < b,
such that
φ0 (c) = 0, φ00 (c) < 0.
By using the change of variable from t to ζ, defined by

φ(t) − φ(c) = −ζ 2 ,

deduce an asymptotic expansion for J(λ) as λ → ∞. Show that the leading-order term
gives
 2π  12
J(λ) ∼ eλφ(c) F (c) .
λ|φ00 (c)|

The gamma function Γ(x) is defined for x > 0 by

Z ∞
Γ(x) = e(x−1) log t−t dt.
0

By means of the substitution t = (x − 1)s, or otherwise, deduce that

1 √  1 
Γ(x + 1) ∼ x(x+ 2 ) e−x 2π 1 + + ...
12x
as x → ∞.

Part II 2004
 31

A2/18 Nonlinear Waves

(i) Let u(x, t) satisfy the Burgers equation

∂u ∂u ∂2u
+u = ν 2,
∂t ∂x ∂x

where ν is a positive constant. Consider solutions of the form u = u(X), where X = x−U t
and U is a constant, such that

∂u ∂u
u → u2 , → 0 as X → −∞; u → u1 , → 0 as X → ∞ ,
∂X ∂X
with u2 > u1 .
Show that U satisfies the so-called shock condition
1
U= (u2 + u1 ).
2
By using the factorisation

1 2 1
u − U u + A = (u − u1 )(u − u2 ),
2 2
where A is the constant of integration, express u in terms of X, u1 , u2 and ν.
(ii) According to shallow-water theory, river waves are characterised by the PDEs

∂v ∂v ∂h v2
+v + g cos α = g sin α − Cf ,
∂t ∂x ∂x h
∂h ∂h ∂v
+v +h = 0,
∂t ∂x ∂x
where h(x, t) denotes the depth of the river, v(x, t) denotes the mean velocity, α is the
constant angle of inclination, and Cf is the constant friction coefficient.
Find the characteristic velocities and the characteristic form of the equations. Find
the Riemann variables and show that if Cf = 0 then the Riemann variables vary linearly
with t on the characteristics.

Part II 2004
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A3/18 Nonlinear Waves

(i) Let Φ+ (t) and Φ− (t) denote the boundary values of functions which are analytic
inside and outside the unit disc centred on the origin, respectively. Let C denote the
boundary of this disc. Suppose that Φ+ (t) and Φ− (t) satisfy the jump condition

Φ+ (t) = t−2 Φ− (t) + t−1 + α(t−1 + t − t−3 ), t ∈ C,

where α is a constant.
Find the canonical solution of the associated homogeneous Riemann-Hilbert prob-
lem. Write down the orthogonality conditions.
(ii) Consider the linear singular integral equation

t − t−1
I
ψ(τ )
(t + t−1 )ψ(t) + dτ = 2 + 2α(1 + t2 − t−2 ), (∗)
πi C τ −t
H
where denotes the principal value integral.
Show that the associated Riemann-Hilbert problem has the jump condition defined
in Part (i) above. Using this fact, find the value of the constant α that allows equation
(∗) to have a solution. For this particular value of α find the unique solution ψ(t).

Part II 2004
 33

A4/23 Nonlinear Waves

Let ψ(k; x, t) satisfy the linear integral equation
Z
3 ψ(l; x, t) 3
ψ(k; x, t) + iei(kx+k t)
dλ(l) = ei(kx+k t) ,
L l+k

where the measure dλ(k) and the contour L are such that ψ(k; x, t) exists and is unique.
Let q(x, t) be defined in terms of ψ(k; x, t) by
Z
∂
q(x, t) = − ψ(k; x, t)dλ(k).
∂x L

(a) Show that

Z
3 (M ψ)(l; x, t)
(M ψ) + iei(kx+k t)
dλ(l) = 0,
L l+k

where
∂2ψ ∂ψ
Mψ ≡ 2
− ik + qψ.
∂x ∂x

(b) Show that

Z Z
i(kx+k3 t) (N ψ)(l; x, t) 3 (M ψ)(l; x, t)
(N ψ) + ie dλ(l) = 3kei(kx+k t) dλ(l),
L l+k L l+k

where
∂ψ ∂ 3 ψ ∂ψ
Nψ ≡ + 3
+ 3q .
∂t ∂x ∂x

(c) By recalling that the KdV equation

∂q ∂3q ∂q
+ 3 + 6q =0
∂t ∂x ∂x
admits the Lax pair
M ψ = 0, N ψ = 0,
write down an expression for dλ(l) which gives rise to the one-soliton solution of the KdV
equation. Write down an expression for ψ(k; x, t) and for q(x, t).

Part II 2004
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A1/1 B1/1 Markov Chains

(i) Give the definitions of a recurrent and a null recurrent irreducible Markov chain.
Let (Xn ) be a recurrent Markov chain with state space I and irreducible transition
matrix P = (pij ). Prove that the vectors γ k = (γjk , j ∈ I), k ∈ I, with entries γkk = 1 and

γik = Ek (# of visits to i before returning to k), i 6= k ,

are P -invariant: X
γjk = γik pij .
i∈I

(ii) Let (Wn ) be the birth and death process on Z+ = {0, 1, 2, . . .} with the following
transition probabilities:
1
pi,i+1 = pi,i−1 = , i ≥ 1
2
p01 = 1 .

By relating (Wn ) to the symmetric simple random walk (Yn ) on Z, or otherwise,

prove that (Wn ) is a recurrent Markov chain. By considering invariant measures, or
otherwise, prove that (Wn ) is null recurrent.
Calculate the vectors γ k = (γik , i ∈ Z+ ) for the chain (Wn ), k ∈ Z+ .
Finally, let W0 = 0 and let N be the number of visits to 1 before returning to 0.
Show that P0 (N = n) = (1/2)n , n ≥ 1.
[You may use properties of the random walk (Yn ) or general facts about Markov
chains without proof but should clearly state them.]

Part II 2004
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A2/1 Markov Chains

(i) Let J be a proper subset of the finite state space I of an irreducible Markov chain
(Xn ), whose transition matrix P is partitioned as

J Jc
 
J A B
P = .
Jc C D

If only visits to states in J are recorded, we see a J-valued Markov chain (X̃n ); show that
its transition matrix is
X
P̃ = A + B Dn C = A + B(I − D)−1 C .
n>0

(ii) Local MP Phil Anderer spends his time in London in the Commons (C), in his flat
(F ), in the bar (B) or with his girlfriend (G). Each hour, he moves from one to another
according to the transition matrix P , though his wife (who knows nothing of his girlfriend)
believes that his movements are governed by transition matrix P W :

C F B G
  C F B
C 1/3 1/3 1/3 0  
C 1/3 1/3 1/3
F 0 1/3 1/3 1/3 
P =  PW = F  1/3 1/3 1/3 
B  1/3 0 1/3 1/3 
B 1/3 1/3 1/3
G 1/3 1/3 0 1/3

The public only sees Phil when he is in J = {C, F, B}; calculate the transition matrix P̃
which they believe controls his movements.
Each time the public Phil moves to a new location, he phones his wife; write down
the transition matrix which governs the sequence of locations from which the public Phil
phones, and calculate its invariant distribution.
Phil’s wife notes down the location of each of his calls, and is getting suspicious
– he is not at his flat often enough. Confronted, Phil swears his fidelity and resolves to
dump his troublesome transition matrix, choosing instead

C F B G
 
C 1/4 1/4 1/2 0
F 1/2 1/4 1/4 0 
P∗ = 
B 0 3/8 1/8 1/2 
G 2/10 1/10 1/10 6/10

Will this deal with his wife’s suspicions? Explain your answer.

Part II 2004
 36

A3/1 B3/1 Markov Chains

(i) Give the definition of the time-reversal of a discrete-time Markov chain (Xn ).
Define a reversible Markov chain and check that every probability distribution satisfying
the detailed balance equations is invariant.
(ii) Customers arrive in a hairdresser’s shop according to a Poisson process of rate
λ > 0. The shop has s hairstylists and N waiting places; each stylist is working (on a
single customer) provided that there is a customer to serve, and any customer arriving
when the shop is full (i.e. the numbers of customers present is N + s) is not admitted
and never returns. Every admitted customer waits in the queue and then is served, in the
first-come-first-served order (say), the service taking an exponential time of rate µ > 0; the
service times of admitted customers are independent. After completing his/her service,
the customer leaves the shop and never returns.
Set up a Markov chain model for the number Xt of customers in the shop at
time t ≥ 0. Assuming λ < sµ, calculate the equilibrium distribution π of this chain
and explain why it is unique. Show that (Xt ) in equilibrium is time-reversible, i.e.
∀ T > 0, (Xt , 0 ≤ t ≤ T ) has the same distribution as (Yt , 0 ≤ t ≤ T ) where Yt = XT −t ,
and X0 ∼ π.

A4/1 Markov Chains

(a) Give three definitions of a continuous-time Markov chain with a given Q-matrix
on a finite state space: (i) in terms of holding times and jump probabilities, (ii) in terms
of transition probabilities over small time intervals, and (iii) in terms of finite-dimensional
distributions.
(b) A flea jumps clockwise on the vertices of a triangle; the holding times are
independent exponential random variables of rate one. Find the eigenvalues of the
corresponding Q-matrix and express transition probabilities pxy (t), t ≥ 0, x, y = A, B, C,
in terms of these roots. Deduce the formulas for the sums
∞ ∞ ∞
X t3n X t3n+1 X t3n+2
S0 (t) = , S1 (t) = , S2 (t) = ,
n=0
(3n)! n=0
(3n + 1)! n=0
(3n + 2)!
√ √
in terms of the functions et , e−t/2 , cos( 3t/2) and sin( 3t/2).
Find the limits
lim e−t Sj (t), j = 0, 1, 2 .
t→∞

What is the connection between the decompositions et = S0 (t) + S1 (t) + S2 (t) and
t
e = cosh t + sinh t?

Part II 2004
 37

A1/2 B1/2 Principles of Dynamics

(i) In Hamiltonian mechanics the action is written
Z  
S= dt pa q̇ a − H(q a , pa , t) . (1)

Starting from Maupertius’ principle δS = 0, derive Hamilton’s equations

∂H ∂H
q̇ a = , ṗa = − .
∂pa ∂q a

Show that H is a constant of the motion if ∂H/∂t = 0. When is pa a constant of the

motion?
(ii) Consider the action S given in Part (i), evaluated on a classical path, as a function
of the final coordinates qfa and final time tf , with the initial coordinates and the initial
time held fixed. Show that S(qfa , tf ) obeys

∂S ∂S
= paf , = −H(qfa , paf , tf ) . (2)
∂qfa ∂tf

Now consider a simple harmonic oscillator with H = 12 (p2 + q 2 ). Setting the initial
time and the initial coordinate to zero, find the classical solution for p and q with final
coordinate q = qf at time t = tf . Hence calculate S(tf , qf ), and explicitly verify (2) in
this case.

Part II 2004
 38

A2/2 B2/1 Principles of Dynamics

(i) Consider a light rigid circular wire of radius a and centre O. The wire lies in
a vertical plane, which rotates about the vertical axis through O. At time t the plane
containing the wire makes an angle φ(t) with a fixed vertical plane. A bead of mass m is
threaded onto the wire. The bead slides without friction along the wire, and its location
is denoted by A. The angle between the line OA and the downward vertical is θ(t).
Show that the Lagrangian of the system is

ma2 2 ma2 2 2
θ̇ + φ̇ sin θ + mga cos θ .
2 2
Calculate two independent constants of the motion, and explain their physical significance.
(ii) A dynamical system has Hamiltonian H(q, p, λ), where λ is a parameter. Consider
an ensemble of identical systems chosen so that the number density of systems, f (q, p, t),
in the phase space element dq dp is either zero or one. Prove Liouville’s Theorem, namely
that the total area of phase space occupied by the ensemble is time-independent.
Now consider a single system undergoing periodic motion q(t), p(t). Give a heuristic
argument based on Liouville’s Theorem to show that the area enclosed by the orbit,
I
I = p dq ,

is approximately conserved as the parameter λ is slowly varied (i.e. that I is an adiabatic

invariant).
Consider H(q, p, λ) = 21 p2 + λq 2n , with n a positive integer. Show that as λ is
slowly varied the energy of the system, E, varies as

E ∝ λ1/(n+1) .

Part II 2004
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A3/2 Principles of Dynamics

(i) Explain the concept of a canonical transformation from coordinates (q a , pa ) to
(Q , P a ). Derive the transformations corresponding to generating functions F1 (t, q a , Qa )
a

and F2 (t, q a , P a ).
(ii) A particle moving in an electromagnetic field is described by the Lagrangian
1  ẋ · A 
L = mẋ2 − e φ − ,
2 c
where c is constant.
(a) Derive the equations of motion in terms of the electric and magnetic fields E
and B.
(b) Show that E and B are invariant under the gauge transformation
1 ∂Λ
A → A + ∇Λ, φ→φ− , (1)
c ∂t
for arbitrary Λ(t, x).
(c) Construct the Hamiltonian. Find the generating function F2 for the canonical
transformation which implements the gauge transformation (1).

A4/2 Principles of Dynamics

Consider a system of coordinates rotating with angular velocity ω relative to an
inertial coordinate system.
Show that if a vector v is changing at a rate dv/dt in the inertial system, then it
is changing at a rate
dv  dv
 = −ω∧v

dt rot dt
with respect to the rotating system.
A solid body rotates with angular velocity ω in the absence of external torque.
Consider the rotating coordinate system aligned with the principal axes of the body.
(a) Show that in this system the motion is described by the Euler equations
dω1 
 dω2  dω3 
I1  = ω2 ω3 (I2 − I3 ) , I2  = ω3 ω1 (I3 − I1 ) , I3  = ω1 ω2 (I1 − I2 ) ,
 
dt rot dt rot dt rot
where (ω1 , ω2 , ω3 ) are the components of the angular velocity in the rotating system and
I1,2,3 are the principal moments of inertia.
(b) Consider a body with three unequal moments of inertia, I3 < I2 < I1 . Show
that rotation about the 1 and 3 axes is stable to small perturbations, but rotation about
the 2 axis is unstable.
(c) Use the Euler equations to show that the kinetic energy, T , and the magnitude
of the angular momentum, L, are constants of the motion. Show further that
2T I3 ≤ L2 ≤ 2T I1 .

Part II 2004
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A1/3 Functional Analysis

(i) Let H be a Hilbert space, and let M be a non-zero closed vector subspace of H.
For x ∈ H, show that there is a unique closest point PM (x) to x in M .
(ii) (a) Let x ∈ H. Show that x − PM (x) ∈ M ⊥ . Show also that if y ∈ M and
x − y ∈ M ⊥ then y = PM (x).
L ⊥
(b) Deduce that H = M M .

(c) Show that the map PM from H to M is a continuous linear map, with ||PM || = 1.

(d) Show that PM is the projection onto M along M ⊥ .

Now suppose that A is a subspace of H that is not necessarily closed. Explain why
A⊥ = {0} implies that A is dense in H.

Give an example of a subspace of l2 that is dense in l2 but is not equal to l2 .

A2/3 B2/2 Functional Analysis

(i) Prove Riesz’s Lemma, that if V is a normed space and A is a vector subspace of
V such that for some 0 6 k < 1 we have d(x, A) 6 k for all x ∈ V with ||x|| = 1, then A
is dense in V . [Here d(x, A) denotes the distance from x to A.]

Deduce that any normed space whose unit ball is compact is finite-dimensional.
[You may assume that every finite-dimensional normed space is complete.]

Give an example of a sequence f1 , f2 , . . . in an infinite-dimensional normed space

such that ||fn || 6 1 for all n, but f1 , f2 , . . . has no convergent subsequence.
(ii) Let V be a vector space, and let ||.||1 and ||.||2 be two norms on V . What does it
mean to say that ||.||1 and ||.||2 are equivalent?

Show that on a finite-dimensional vector space all norms are equivalent. Deduce
that every finite-dimensional normed space is complete.

Exhibit two norms on the vector space l1 that are not equivalent.

In addition, exhibit two norms on the vector space l∞ that are not equivalent.

Part II 2004
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A3/3 B3/2 Functional Analysis

(i) Let H be an infinite-dimensional Hilbert space. Show that H has a (countable)
orthonormal basis if and only if H has a countable dense subset. [You may assume
familiarity with the Gram-Schmidt process.]

State and prove Bessel’s inequality.

(ii) State Parseval’s equation. Using this, prove that if H has a countable dense subset
then there is a surjective isometry from H to l2 .

Explain carefully why the functions einθ , n ∈ Z, form an orthonormal basis for
L2 (T).

A4/3 Functional Analysis

State and prove the Dominated Convergence Theorem. [You may assume the
Monotone Convergence Theorem.]

Let a and p be real numbers, with a > 0. Prove carefully that

Z ∞
p
e−ax sin px dx = .
0 a2 + p2

[Any standard results that you use should be stated precisely.]

Part II 2004
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A1/4 B1/3 Groups, Rings and Fields

(i) Let R be a commutative ring. Define the terms prime ideal and maximal ideal,
and show that if an ideal M in R is maximal then M is also prime.
(ii) Let P be a non-trivial prime ideal in the commutative ring R (‘non-trivial’ meaning
that P 6= {0} and P 6= R). If P has finite index as a subgroup of R, show that P is also
maximal. Give an example to show that this may fail, if the assumption of finite index is
omitted. Finally, show that if R is a principal ideal domain, then every non-trivial prime
ideal in R is maximal.

A2/4 B2/3 Groups, Rings and Fields

(i) State Gauss’ Lemma on polynomial irreducibility. State and prove Eisenstein’s
criterion.
(ii) Which of the following polynomials are irreducible over Q? Justify your answers.
(a) x7 − 3x3 + 18x + 12
(b) x4 − 4x3 + 11x2 − 3x − 5
(c) 1 + x + x2 + . . . + xp−1 with p prime
[Hint: consider substituting y = x − 1.]
(d) xn + px + p2 with p prime.
[Hint: show any factor has degree at least two, and consider powers of p dividing
coefficients.]

A3/4 Groups, Rings and Fields

(i) Let K 6 C be a field and L 6 C a finite normal extension of K. If H is a finite
subgroup of order m in the Galois group G(L | K), show that L is a normal extension of
the H-invariant subfield I(H) of degree m and that G(L | I(H)) = H. [You may assume
the theorem of the primitive element.]
√
(ii) Show that the splitting field over Q of the polynomial x4 +2 is Q[ 4 2, i] and deduce
that its Galois group has order 8. Exhibit a subgroup of order 4 of the Galois group, and
determine the corresponding invariant subfield.

Part II 2004
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A4/4 Groups, Rings and Fields

(a) Let t be the maximal power of the prime p dividing the order of the finite group
G, and let N (pt ) denote the number of subgroups of G of order pt . State clearly the
numerical restrictions on N (pt ) given by the Sylow theorems.
If H and K are subgroups of G of orders r and s respectively, and their intersection
H ∩ K has order t, show the set HK = {hk : h ∈ H, k ∈ K} contains rs/t elements.
(b) The finite group G has 48 elements. By computing the possible values of N (16),
show that G cannot be simple.

Part II 2004
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A1/5 B1/4 Electromagnetism

(i) Show that the work done in assembling a localised charge distribution ρ(r) in a
region V with an associated potential φ(r) is

1
R
W = 2 V
ρ(r)φ(r) dτ,

and that this can be written as an integral over all space

W = 12 0 |E|2 dτ,
R

where the electric field E = −∇φ.

(ii) What is the force per unit area on an infinite plane conducting sheet with a charge
density σ per unit area (a) if it is isolated in space and (b) if the electric field vanishes on
one side of the sheet?
An infinite cylindrical capacitor consists of two concentric cylindrical conductors
with radii a, b (a < b), carrying charges ±q per unit length respectively. Calculate the
capacitance per unit length and the energy per unit length. Next determine the total
force on each conductor, and calculate the rate of change of energy of the inner and outer
conductors if they are moved radially inwards and outwards respectively with speed v.
What is the corresponding rate of change of the capacitance?

A2/5 Electromagnetism
(i) Write down the general solution of Poisson’s equation. Derive from Maxwell’s
equations the Biot-Savart law for the magnetic field of a steady localised current distribu-
tion.
(ii) A plane rectangular loop with sides of length a and b lies in the plane z = 0 and
is centred on the origin. Show that when r = |r|  a, b, the vector potential A(r) is given
approximately by
µ0 m ∧ r
A(r) = ,
4π r3
where m = Iabẑ is the magnetic moment of the loop.
Hence show that the magnetic field B(r) at a great distance from an arbitrary small
plane loop of area A, situated in the xy-plane near the origin and carrying a current I, is
given by
µ0 IA
3xz, 3yz, 2r2 − 3x2 − 3y 2 .


B(r) = 5
4πr

Part II 2004
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A3/5 B3/3 Electromagnetism

(i) State Maxwell’s equations and show that the electric field E and the magnetic
field B can be expressed in terms of a scalar potential φ and a vector potential A. Hence
derive the inhomogeneous wave equations that are satisfied by φ and A respectively.
(ii) The plane x = 0 separates a vacuum in the half-space x < 0 from a perfectly
conducting medium occupying the half-space x > 0. Derive the boundary conditions on
E and B at x = 0.
A plane electromagnetic wave with a magnetic field B = B(t, x, z)ŷ, travelling in
the xz-plane at an angle θ to the x-direction, is incident on the interface at x = 0. If the
wave has frequency ω show that the total magnetic field is given by
h
i
ωx ωz


B = B0 cos c cos θ exp i c sin θ − ωt ŷ,

where B0 is a constant. Hence find the corresponding electric field E, and obtain the
surface charge density and the surface current at the interface.

A4/5 Electromagnetism
Consider a frame S 0 moving with velocity v relative to the laboratory frame S
where |v|2  c2 . The electric and magnetic fields in S are E and B, while those measured
in S 0 are E0 and B0 . Given that B0 = B, show that
I I
0
E · dl = (E + v ∧ B) · dl,
Γ Γ

for any closed circuit Γ and hence that E0 = E + v ∧ B.

Now consider a fluid with electrical conductivity σ and moving with velocity v(r).
Use Ohm’s law in the moving frame to relate the current density j to the electric field E
in the laboratory frame, and show that if j remains finite in the limit σ → ∞ then

∂B
= ∇ ∧ (v ∧ B).
∂t
R
The magnetic helicity H in a volume V is given by V A · B dτ where A is the
vector potential. Show that if the normal components of v and B both vanish on the
surface bounding V then dH/dt = 0.

Part II 2004
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B2/4 Dynamics of Differential Equations

(i) Define carefully what is meant by a Hopf bifurcation in a two-dimensional dynam-
ical system. Write down the normal form for this bifurcation, correct to cubic order, and
distinguish between bifurcations of supercritical and subcritical type. Describe, without
detailed calculations, how a general two-dimensional system with a Hopf bifurcation at
the origin can be reduced to normal form by a near-identity transformation.
(ii) A Takens-Bogdanov bifurcation of a fixed point of a two-dimensional system is
characterised by a Jacobian with the canonical form
 
0 1
A=
0 0

at the bifurcation point. Consider the system

ẋ = y + α1 x2 + β1 xy + γ1 y 2
ẏ = α2 x2 + β2 xy + γ2 y 2 .

Show that a near-identity transformation of the form

ξ = x + a1 x2 + b1 xy + c1 y 2
η = y + a2 x2 + b2 xy + c2 y 2

exists that reduces the system to the normal (canonical) form, correct up to quadratic
terms,
ξ˙ = η, η̇ = α2 ξ 2 + (β2 + 2α1 )ξη.

It is known that the general form of the equations near the bifurcation point can
be written (setting p = α2 , q = β2 + 2α1 )

ξ˙ = η, η̇ = λξ + µη + pξ 2 + qξη.

Find all the fixed points of this system, and the values of λ, µ for which these fixed
points have (a) steady state bifurcations and (b) Hopf bifurcations.

Part II 2004
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B3/4 Dynamics of Differential Equations

(i) Describe the use of the stroboscopic method for obtaining approximate solutions to
the second order equation
ẍ + x = f (x, ẋ, t)
when ||  1. In particular, by writing x = R cos(t + φ), ẋ = −R sin(t + φ), obtain
expressions in terms of f for the rate of change of R and φ. Evaluate these expressions
when f = x2 cos t.
(ii) In planetary orbit theory a crude model of an orbit subject to perturbation from
a distant body is given by the equation

d2 u
+ u = λ − δ 2 u−2 − 2δ 3 u−3 cos θ,
dθ2

where 0 < δ  1, (u−1 , θ) are polar coordinates in the plane, and λ is a positive constant.
(a) Show that when δ = 0 all bounded orbits are closed.
(b) Now suppose δ 6= 0, and look for almost circular orbits with u = λ+δw(θ)+aδ 2 ,
where a is a constant. By writing w = R(θ) cos(θ + φ(θ)), and by making a suitable choice
of the constant a, use the stroboscopic method to find equations for dw/dθ and dφ/dθ.
By writing z = R exp(iφ) and considering dz/dθ, or otherwise, determine R(θ) and φ(θ)
in the case R(0) = R0 , φ(0) = 0. Hence describe the orbits of the system.

Part II 2004
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A1/7 B1/12 Logic, Computation and Set Theory

(i) State and prove the Knaster-Tarski Fixed-Point Theorem.
(ii) A subset S of a poset X is called an up-set if whenever x, y ∈ X satisfy x ∈ S
and x 6 y then also y ∈ S. Show that the set of up-sets of X (ordered by inclusion) is a
complete poset.

Let X and Y be totally ordered sets, such that X is isomorphic to an up-set in Y

and Y is isomorphic to the complement of an up-set in X. Prove that X is isomorphic to
Y . Indicate clearly where in your argument you have made use of the fact that X and Y
are total orders, rather than just partial orders.

[Recall that posets X and Y are called isomorphic if there exists a bijection f from
X to Y such that, for any x, y ∈ X, we have f (x) 6 f (y) if and only if x 6 y.]

B2/11 Logic, Computation and Set Theory

Define the sets Vα , α ∈ ON . Show that each Vα is transitive, and explain why
Vα ⊆ Vβ whenever α 6 β. Prove that every set x is a member of some Vα .

Which of the following are true and which are false? Give proofs or counterexamples
as appropriate. You may assume standard properties of rank.

(a) If the rank of a set x is a (non-zero) limit then x is infinite.

(b) If the rank of a set x is a successor then x is finite.

(c) If the rank of a set x is countable then x is countable.

A3/8 B3/11 Logic, Computation and Set Theory

(i) State and prove the Compactness Theorem for first-order predicate logic.

State and prove the Upward Löwenheim-Skolem Theorem.

[You may use the Completeness Theorem for first-order predicate logic.]
(ii) For each of the following theories, either give axioms (in the language of posets)
for the theory or prove carefully that the theory is not axiomatisable.

(a) The theory of posets having no maximal element.

(b) The theory of posets having a unique maximal element.

(c) The theory of posets having infinitely many maximal elements.

(d) The theory of posets having finitely many maximal elements.

(e) The theory of countable posets having a unique maximal element.

Part II 2004
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A4/8 B4/10 Logic, Computation and Set Theory

Write an essay on recursive functions. Your essay should include a sketch of
why every computable function is recursive, and an explanation of the existence of a
universal recursive function, as well as brief discussions of the Halting Problem and of the
relationship between recursive sets and recursively enumerable sets.
[You may assume that every recursive function is computable. You do not need to
give proofs that particular functions to do with prime-power decompositions are recursive.]

Part II 2004
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A1/12 B1/15 Principles of Statistics

(i) What does it mean to say that a family {f (·|θ) : θ ∈ Θ} of densities is an
exponential family?
Consider the family of densities on (0, ∞) parametrised by the positive parameters
a, b and defined by

a exp(−(a − bx)2 /2x)

f (x|a, b) = √ (x > 0).
2πx3

Prove that this family is an exponential family, and identify the natural parameters and
the reference measure.
(ii) Let (X1 , . . . , Xn ) be a sample drawn from the above distribution. Find the
maximum-likelihood estimators of the parameters (a, b). Find the Fisher information
matrix of the family (in terms of the natural parameters). Briefly explain the significance
of the Fisher information matrix in relation to unbiased estimation. Compute the mean
of X1 and of X1−1 .

Part II 2004
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A2/11 B2/16 Principles of Statistics

(i) In the context of a decision-theoretic approach to statistics, what is a loss function?
a decision rule? the risk function of a decision rule? the Bayes risk of a decision rule? the
Bayes rule with respect to a given prior distribution?
Show how the Bayes rule with respect to a given prior distribution is computed.
(ii) A sample of n people is to be tested for the presence of a certain condition. A
single real-valued observation is made on each one; this observation comes from density
f0 if the condition is absent, and from density f1 if the condition is present. Suppose
θi = 0 if the ith person does not have the condition, θi = 1 otherwise, and suppose that
the prior distribution for the θi is that they are independent with common distribution
P (θi = 1) = p ∈ (0, 1), where p is known. If Xi denotes the observation made on the ith
person, what is the posterior distribution of the θi ?
Now suppose that the loss function is defined by
n
X
L0 (θ, a) ≡ (αaj (1 − θj ) + β(1 − aj )θj )
j=1

for action a ∈ [0, 1]n , where α, β are positive constants. If πj denotes the posterior
probability that θj = 1 given the data, prove that the Bayes rule for this prior and this
loss function is to take aj = 1 if πj exceeds the threshold value α/(α + β), and otherwise
to take aj = 0.
In an attempt to control the proportion of false positives, it is proposed to use a
different loss function, namely,
 P 
θj aj
L1 (θ, a) ≡ L0 (θ, a) + γI{ aj >0} 1 − P
P ,
aj

where γ > 0. Prove that the Bayes rule is once again a threshold rule, that is, we take
action aj = 1 if and only if πj > λ, and determine λ as fully as you can.

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A3/12 B3/15 Principles of Statistics

(i) What is a sufficient statistic? What is a minimal sufficient statistic? Explain the
terms nuisance parameter and ancillary statistic.
(ii) Let U1 , . . . , Un be independent random variables with common uniform([0, 1])
distribution, and suppose you observe Xi ≡ aUi−β , i = 1, . . . , n, where the positive
parameters a, β are unknown. Write down the joint density of X1 , . . . , Xn and prove
that the statistic
n
Y
(m, p) ≡ ( min {Xj }, Xj )
16j6n
j=1

is minimal sufficient for (a, β). Find the maximum-likelihood estimator (â, β̂) of (a, β).
Regarding β as the parameter of interest and a as the nuisance parameter, is m
ancillary? Find the mean and variance of β̂. Hence find an unbiased estimator of β.

A4/13 B4/15 Principles of Statistics

Suppose that θ ∈ Rd is the parameter of a non-degenerate exponential family.
Derive the asymptotic distribution of the maximum-likelihood estimator θ̂n of θ based
on a sample of size n. [You may assume that the density is infinitely differentiable with
respect to the parameter, and that differentiation with respect to the parameter commutes
with integration.]

Part II 2004
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A1/11 B1/16 Stochastic Financial Models

(i) What does it mean to say that U is a utility function? What is a utility function
with constant absolute risk aversion (CARA)?
Let St ≡ (St1 , . . . , Std )T denote the prices at time t = 0, 1 of d risky assets, and
suppose that there is also a riskless zeroth asset, whose price at time 0 is 1, and whose price
at time 1 is 1 + r. Suppose that S1 has a multivariate Gaussian distribution, with mean µ1
and non-singular covariance V . An agent chooses at time 0 a portfolio θ = (θ1 , . . . , θd )T
of holdings of the d risky assets, at total cost θ · S0 , and at time 1 realises his gain
X = θ · ( S1 − (1 + r)S0 ). Given that he wishes the mean of X to be equal to m, find the
smallest value that the variance v of X can be. What is the portfolio that achieves this
smallest variance? Hence sketch the region in the (v, m) plane of pairs (v, m) that can be
achieved by some choice of θ, and indicate the mean-variance efficient frontier.
(ii) Suppose that the agent has a CARA utility with coefficient γ of absolute risk
aversion. What portfolio will he choose in order to maximise EU (X)? What then is the
mean of X?
Regulation requires that the agent’s choice of portfolio θ has to satisfy the value-
at-risk (VaR) constraint √
m > −L + a v,
where L > 0 and a > 0 are determined by the regulatory p authority. Show that this
constraint has no effect on the agent’s decision if κ ≡ µ · V −1 µ > a. If κ < a, will this
constraint necessarily affect the agent’s choice of portfolio?

A3/11 B3/16 Stochastic Financial Models

(i) Consider a single-period binomial model of a riskless asset (asset 0), worth 1 at
time 0 and 1 + r at time 1, and a risky asset (asset 1), worth 1 at time 0 and worth u at
time 1 if the period was good, otherwise worth d. Assuming that

d<1+r <u (∗)

show how any contingent claim Y to be paid at time 1 can be priced and exactly replicated.
Briefly explain the significance of the condition (∗), and indicate how the analysis of the
single-period model extends to many periods.
(ii) Now suppose that u = 5/3, d = 2/3, r = 1/3, and that the risky asset is worth
S0 = 864 = 25 × 33 at time zero. Show that the time-0 value of an American put option
with strike K = S0 and expiry at time t = 3 is equal to 79, and find the optimal exercise
policy.

Part II 2004
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A4/12 B4/16 Stochastic Financial Models

What is Brownian motion (Bt )t>0 ? Briefly explain how Brownian motion can be
considered as a limit of simple random walks. State the Reflection Principle for Brownian
motion, and use it to derive the distribution of the first passage time τa ≡ inf{t : Bt = a}
to some level a > 0.
Suppose that Xt = Bt + ct, where c > 0 is constant. Stating clearly any results to
(c)
which you appeal, derive the distribution of the first-passage time τa ≡ inf{t : Xt = a}
to a > 0.
Now let σa ≡ sup{t : Xt = a}, where a > 0. Find the density of σa .

Part II 2004
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A2/13 B2/22 Foundations of Quantum Mechanics

(i) The creation and annihilation operators for a harmonic oscillator of angular
frequency ω satisfy the commutation relation [a, a† ] = 1. Write down an expression for
the Hamiltonian H in terms of a and a† .
There exists a unique ground state |0i of H such that a|0i = 0. Explain how the
space of eigenstates |ni, n = 0, 1, 2, . . . of H is formed, and deduce the eigenenergies for
these states. Show that
√ √
a|ni = n|n − 1i , a† |ni = n + 1|n + 1i .

(ii) Write down the number operator N of the harmonic oscillator in terms of a and
a† . Show that
N |ni = n|ni .
The operator Kr is defined to be

a†r ar
Kr = , r = 0, 1, 2, . . . .
r!
Show that Kr commutes with N . Show also that
n!
(n−r)! r! |ni r≤n,
(
Kr |ni =
0 r>n.

By considering the action of Kr on the state |ni show that

∞
X
(−1)r Kr = |0ih0| .
r=0

Part II 2004
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A3/13 B3/21 Foundations of Quantum Mechanics

(i) A quantum mechanical system consists of two identical non-interacting particles
with associated single-particle wave functions ψi (x) and energies Ei , i = 1, 2, . . . ., where
E1 < E2 < . . . . Show how the states for the two lowest energy levels of the system are
constructed and discuss their degeneracy when the particles have (a) spin 0, (b) spin 1/2.
(ii) The Pauli matrices are defined to be
     
0 1 0 −i 1 0
σ1 = , σ2 = , σ3 = .
1 0 i 0 0 −1

State how the spin operators s1 , s2 , s3 may be expressed in terms of the Pauli matrices,
and show that they describe states with total angular momentum 12 ~.
An electron is at rest in the presence of a magnetic field B = (B, 0, 0), and
experiences an interaction potential −µσ · B. At t = 0 the state of the electron is the
eigenstate of s3 with eigenvalue 12 ~. Calculate the probability that at later time t the
electron will be measured to be in the eigenstate of s3 with eigenvalue 21 ~.

Part II 2004
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A4/15 B4/22 Foundations of Quantum Mechanics

The states of the hydrogen atom are denoted by |nlmi with l < n, −l ≤ m ≤ l and
associated energy eigenvalue En , where

e2
En = − .
8π0 a0 n2

A hydrogen atom is placed in a weak electric field with interaction Hamiltonian

H1 = −eEz .

a) Derive the necessary perturbation theory to show that to O(E 2 ) the change in the
energy associated with the state |100i is given by
 2
∞ n−1 l h100|z|nlmi
 
X X X
∆E1 = e2 E 2 . (∗)
n=2 l=0 m=−l
E1 − En

The wavefunction of the ground state |100i is

1
ψn=1 (r) = e−r/a0 .
(πa30 )1/2

By replacing En , ∀ n > 1, in the denominator of (∗) by E2 show that

32π
|∆E1 | < 0 E 2 a30 .
3

b) Find a matrix whose eigenvalues are the perturbed energies to O(E) for the states
|200i and |210i. Hence, determine these perturbed energies to O(E) in terms of the
matrix elements of z between these states.
[Hint:
hnlm|z|nlmi = 0 ∀ n, l, m
0 0
hnlm|z|nl m i = 0 ∀ n, l, l0 , m, m0 , m 6= m0
]

Part II 2004
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A1/15 B1/24 General Relativity

(i) What is an affine parameter λ of a timelike or null geodesic? Prove that for a
timelike geodesic one may take λ to be proper time τ . The metric

ds2 = −dt2 + a2 (t) dx2 ,

with ȧ(t) > 0 represents an expanding universe. Calculate the Christoffel symbols.
(ii) Obtain the law of spatial momentum conservation for a particle of rest mass m in
the form
dx
ma2 = p = constant.
dτ
Assuming that the energy E = m dt/dτ , derive an expression for E in terms of m, p and
a(t) and show that the energy is not conserved but rather that it decreases with time. In
particular, show that if the particle is moving extremely relativistically then the energy
decreases as a−1 (t), and if it is moving non-relativistically then the kinetic energy, E − m,
decreases as a−2 (t).
Show that the frequency ωe of a photon emitted at time te will be observed at time
to to have frequency
a(te )
ωo = ωe .
a(to )

A2/15 B2/24 General Relativity

(i) State and prove Birkhoff’s theorem.
(ii) Derive the Schwarzschild metric and discuss its relevance to the problem of
gravitational collapse and the formation of black holes.
[Hint: You may assume that the metric takes the form

ds2 = −eν(r,t) dt2 + eλ(r,t) dr2 + r2 (dθ2 + sin2 θ dφ2 ),

and that the non-vanishing components of the Einstein tensor are given by

e2ν+λ λ̇ eλ
Gtt = (−1 + eλ + rλ0 ), Grt = e(ν+λ)/2 , Grr = (1 − e−λ + rν 0 ),
r2 r r2
h 2 i
Gθθ = 14 r2 e−λ 2ν 00 + (ν 0 )2 + (ν 0 − λ0 ) − ν 0 λ0 − 14 r2 e−ν 2λ̈ + (λ̇)2 − λ̇ν̇ ,
 
r
Gtr = Grt and Gφφ = sin2 θ Gθθ .]

Part II 2004
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A4/17 B4/25 General Relativity

Starting from the Ricci identity

Va;b;c − Va;c;b = Ve Re abc ,

give an expression for the curvature tensor Re abc of the Levi-Civita connection in terms
of the Christoffel symbols and their partial derivatives. Using local inertial coordinates,
or otherwise, establish that

Re abc + Re bca + Re cab = 0. (∗)

A vector field with components V a satisfies

Va;b + Vb;a = 0. (∗∗)

Show, using equation (∗) that

Va;b;c = Ve Re cba ,
and hence that
Va;b ;b + Ra c Vc = 0,
where Rab is the Ricci tensor. Show that equation (∗∗) may be written as

(∂c gab )V c + gcb ∂a V c + gac ∂b V c = 0. (∗∗∗)

If the metric is taken to be the Schwarzschild metric

2M −1
ds2 = − 1 − 2M
dr2 + r2 (dθ2 + sin2 θ dφ2 ),



r dt2 + 1 − r

show that V a = δ a 0 is a solution of (∗∗∗). Calculate V a ;a .

Electromagnetism can be described by a vector potential Aa and a Maxwell field
tensor Fab satisfying

Fab = Ab;a − Aa;b and Fab ;b = 0. (∗∗∗∗)

The divergence of Aa is arbitrary and we may choose Aa ;a = 0. With this choice show
that in a general spacetime
Aa;b ;b − Ra c Ac = 0.
Hence show that in the Schwarzschild spacetime a tensor field whose only non-trivial
components are Ftr = −Frt = Q/r2 , where Q is a constant, satisfies the field equations
(∗∗∗∗).

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A1/20 B1/20 Numerical Analysis

(i) Define the Backward Difference Formula (BDF) method for ordinary differential
equations and derive its two-step version.
(ii) Prove that the interval (−∞, 0) belongs to the linear stability domain of the two-
step BDF method.

A2/19 B2/20 Numerical Analysis

(i) The five-point equations, which are obtained when the Poisson equation ∇2 u = f
(with Dirichlet boundary conditions) is discretized in a square, are

−um−1,n − um,n−1 − um+1,n − um,n+1 + 4um,n = fm,n , m, n = 1, 2, . . . , M,

where u0,n , uM +1,n , um,0 , um,M +1 = 0 for all m, n = 1, 2, . . . , M .

Formulate the Gauss–Seidel method for the above linear system and prove its
convergence. In the proof you should carefully state any theorems you use. [You may
use Part (ii) of this question.]
(ii) By arranging the two-dimensional arrays {um,n }m,n=1,...,M and {bm,n }m,n=1,...,M
2 2
into the column vectors u ∈ RM and b ∈ RM respectively, the linear system described
in Part (i) takes the matrix form Au = b. Prove that, regardless of the ordering of the
points on the grid, the matrix A is symmetric and positive definite.

A3/19 B3/20 Numerical Analysis

(i) The diffusion equation

∂u ∂2u
= , 0 6 x 6 1, t > 0,
∂t ∂x2
with the initial condition u(x, 0) = φ(x), 0 6 x 6 1, and with zero boundary conditions at
x = 0 and x = 1, can be solved by the method

un+1
m = unm + µ(unm−1 − 2unm + unm+1 ), m = 1, 2, . . . , M, n > 0,
1
where ∆x = 1/(M +1), µ = ∆t/(∆x)2 , and unm ≈ u(m∆x, n∆t). Prove that µ 6 2 implies
convergence.
(ii) By discretizing the same equation and employing the same notation as in Part (i),
determine conditions on µ > 0 such that the method
1 1  5  1 1 
− µ un+1
m−1 + + µ un+1
m + − µ un+1
m+1 =
12 2 6 12 2
1 1   5   1 1 
+ µ unm−1 + − µ unm + + µ unm+1
12 2 6 12 2
is stable.

Part II 2004
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A4/22 B4/20 Numerical Analysis

Write an essay on the method of conjugate gradients. You should define the method,
list its main properties and sketch the relevant proof. You should also prove that (in
exact arithmetic) the method terminates in a finite number of steps, briefly mention the
connection with Krylov subspaces, and describe the approach of preconditioned conjugate
gradients.

Part II 2004
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A1/6 B1/17 Nonlinear Dynamical Systems

(i) State Liapunov’s First Theorem and La Salle’s Invariance Principle. Use these
results to show that the system

ẍ + k ẋ + sin x = 0, k > 0

has an asymptotically stable fixed point at the origin.

(ii) Define the basin of attraction of an invariant set of a dynamical system.
Consider the equations

ẋ = −x + βxy 2 + x3 , ẏ = −y + βyx2 + y 3 , β > 2.

(a) Find the fixed points of the system and determine their type.
(b) Show that the basin of attraction of the origin includes the union over α of the
regions
4α2 (1 + α2 )(β − 1)
x2 + α2 y 2 < 2 .
β (1 + α2 )2 − 4α2
Sketch these regions for α2 = 1, 1/2, 2 in the case β = 3.

A2/6 B2/17 Nonlinear Dynamical Systems

(i) A linear system in R2 takes the form ẋ = Ax. Explain (without detailed calculation
but by giving examples) how to classify the dynamics of the system in terms of the
determinant and the trace of A. Show your classification graphically, and describe the
dynamics that occurs on the boundaries of the different regions on your diagram.
(ii) A nonlinear system in R2 has the form ẋ = f (x), f (0) = 0. The Jacobian
(linearization) A of f at the origin is non-hyperbolic, with one eigenvalue of A in the
left-hand half-plane. Define the centre manifold for this system, and explain (stating
carefully any results you use) how the dynamics near the origin may be reduced to a
one-dimensional system on the centre manifold.
A dynamical system of this type has the form

ẋ = ax3 + bxy + cx5 + dx3 y + exy 2 + f x7 + gx5 y

ẏ = −y + x2 − x4

Find the coefficients for the expansion of the centre manifold correct up to and
including terms of order x6 , and write down in terms of these coefficients the equation for
the dynamics on the centre manifold up to order x7 . Using this reduced equation, give a
complete set of conditions on the coefficients a, b, c, . . . that guarantee that the origin is
stable.

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A4/6 B4/17 Nonlinear Dynamical Systems

(a) Consider the map G1 (x) = f (x+a), defined on 0 6 x < 1, where f (x) = x [mod 1],
0 6 f < 1, and the constant a satisfies 0 6 a < 1. Give, with reasons, the values of a (if
any) for which the map has (i) a fixed point, (ii) a cycle of least period n, (iii) an aperiodic
orbit. Does the map exhibit sensitive dependence on initial conditions?
Show (graphically if you wish) that if the map has an n-cycle then it has an infinite
number of such cycles. Is this still true if G1 is replaced by f (cx + a), 0 < c < 1?
(b) Consider the map

b
G2 (x) = f (x + a + sin 2πx),
2π

where f (x) and a are defined as in Part (a), and b > 0 is a parameter.
Find the regions of the (a, b) plane for which the map has (i) no fixed points,
(ii) exactly two fixed points.
Now consider the possible existence of a 2-cycle of the map G2 when b  1, and
suppose the elements of the cycle are X, Y with X < 12 . By expanding X, Y, a in powers
of b, so that X = X0 + bX1 + b2 X2 + O(b3 ), and similarly for Y and a, show that

1 b2
a= + sin 4πX0 + O(b3 ).
2 8π

Use this result to sketch the region of the (a, b) plane in which 2-cycles exist. How many
distinct cycles are there for each value of a in this region?

Part II 2004
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A3/6 B3/17 Nonlinear Dynamical Systems

(i) Consider a system in R2 that is almost Hamiltonian:

∂H ∂H
ẋ = + g1 (x, y), ẏ = − + g2 (x, y) ,
∂y ∂x

where
H H = H(x, y) and ||  1. Show that if the system has a periodic orbit C then
g
C 2
dx − g1 dy = 0, and explain how to evaluate this orbit approximately for small .
Illustrate your method by means of the system

ẋ = y + x(1 − x2 ), ẏ = −x.

(ii) Consider the system

ẋ = y, ẏ = x − x3 + y(1 − αx2 ).

(a) Show that when  = 0 the system is Hamiltonian, and find the Hamiltonian.
Sketch the trajectories in the case  = 0. Identify the value Hc of H for which there is a
homoclinic orbit.
(b) Suppose  > 0. Show that the small change ∆H in H around an orbit of the
Hamiltonian system can be expressed to leading order as an integral of the form
Z x2
F(x, H)dx,
x1

where x1 , x2 are the extrema of the x-coordinates of the orbits of the Hamiltonian system,
distinguishing between the cases H < Hc , H > Hc .
(c) Find the value of α, correct to leading order in , at which the system has a
homoclinic orbit.
(d) By examining the eigenvalues of the Jacobian at the origin, determine the
stability of the homoclinic orbit, being careful to state clearly any standard results that
you use.

Part II 2004
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B1/5 Combinatorics
State and prove Menger’s theorem (vertex form).

Let G be a graph of connectivity κ(G) ≥ k and let S, T be disjoint subsets of V (G)

with |S|, |T | ≥ k. Show that there exist k vertex disjoint paths from S to T .

The graph H is said to be k-linked if, for every sequence s1 , . . . , sk , t1 , . . . , tk of 2k

distinct vertices, there exist si – ti paths, 1 ≤ i ≤ k, that are vertex disjoint. By removing
an edge from K2k , or otherwise, show that, for k > 2, H need not be k-linked even if
κ(H) ≥ 2k − 2.

Prove that if |H| = n and δ(H) ≥ 21 (n + 3k) − 2 then H is k-linked.

B2/5 Combinatorics
State and prove Sperner’s lemma on antichains.

The family A ⊂ P[n] is said to split [n] if, for all distinct i, j
∈ [n], there exists
a
A ∈ A with i ∈ A but j ∈
/ A. Prove that if A splits [n] then n ≤ ba/2c , where a = |A|.

Show moreover that,
if A splits [n] and no element of [n] is in more than k < ba/2c
members of A, then n ≤ ka .

B4/1 Combinatorics
Write an essay on Ramsey’s theorem. You should include the finite and infinite
versions, together with some discussion of bounds in the finite case, and give at least one
application.

Part II 2004
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B1/6 Representation Theory

(a) Show that every irreducible complex representation of an abelian group is one-
dimensional.
(b) Show, by example, that the analogue of (a) fails for real representations.
(c) Let the cyclic group of order n act on Cn by cyclic permutation of the standard
basis vectors. Decompose this representation explicitly into irreducibles.

B2/6 Representation Theory

Let H be a group with three generators c, g, h and relations cp = g p = hp = 1,
cg = gc, ch = hc and gh = chg where p is a prime number.
(a) Show that |H| = p3 . Show that the conjugacy classes of H are the singletons
{1}, {c}, . . . , {cp−1 } and the sets {g m hn , cg m hn , . . . , cp−1 g m hn }, as m, n range from
0 to p − 1, but (m, n) 6= (0, 0).
(b) Find p2 1-dimensional representations of H.
(c) Let ω 6= 1 be a pth root of unity. Show that the following defines an irreducible
representation of H on Cp :

ρ(c) = ωId,
ρ(g)ek = ω k ek ,
ρ(h)ep = e1 and ρ(h)ek = ek+1 if k < p

where the ek are the standard basis vectors of Cp .

(d) Show that (b) and (c) cover all irreducible isomorphism classes.

Part II 2004
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B3/5 Representation Theory

Compute the character table for the group A5 of even permutations of five elements.
You may wish to follow the steps below.
(a) List the conjugacy classes in A5 and their orders.
(b) A5 acts on C5 by permuting the standard basis vectors. Show that C5 splits as
C ⊕ V , where C is the trivial 1-dimensional representation and V is irreducible.
(c) By using the formula for the character of the symmetric square S 2 V ,

1
χV (g)2 + χV (g 2 ) ,

χS 2 V (g) =
2

decompose S 2 V to produce a 5-dimensional, irreducible representation, and find

its character.
(d) Show that the exterior square Λ2 V decomposes into two distinct irreducibles and
compute their characters, to complete the character table of A5 .
[Hint: You can save yourself some computational effort if you can explain why the
automorphism of A5 , defined by conjugation by a transposition in S5 , must swap the two
summands of Λ2 V .]

B4/2 Representation Theory

Write an essay on the finite-dimensional representations of SU2 , including a proof
of their complete reducibility, and a description of the irreducible representations and the
decomposition of their tensor products.

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B1/7 Galois Theory

Let L/K be a finite extension of fields. Define the trace TrL/K (x) and norm
NL/K (x) of an element x ∈ L.
Assume now that the extension L/K is Galois, with cyclic Galois group of prime
order p, generated by σ.
Pp−1
i) Show that TrL/K (x) = n=0 σ n (x).
ii) Show that {σ(x) − x | x ∈ L} is a K-vector subspace of L of dimension p − 1.
Deduce that if y ∈ L, then TrL/K (y) = 0 if and only if y = σ(x) − x for some x ∈ L.
[You may assume without proof that TrL/K is surjective for any finite separable extension
L/K.]
iii) Suppose that L has characteristic p. Deduce from (i) that every element of K
can be written as σ(x) − x for some x ∈ L. Show also that if σ(x) = x + 1, then xp − x
belongs to K. Deduce that L is the splitting field over K of X p − X − a for some a ∈ K.

B3/6 Galois Theory

Let K be a field, and G a finite subgroup of K ∗ . Show that G is cyclic.
Define the cyclotomic polynomials Φm , and show from your definition that
Y
Xm − 1 = Φd (X).
d|m

Deduce that Φm is a polynomial with integer coefficients.

Let p be a prime with (m, p) = 1. Let Φm ≡ f1 . . . fr (mod p), where fi ∈ Fp [X]
are irreducible. Show that for each i the degree of fi is equal to the order of p in the group
(Z/mZ)∗ .
Use this to write down an irreducible polynomial of degree 10 over F2 .

B4/3 Galois Theory

Let M/K be a finite Galois extension of fields. Explain what is meant by the Galois
correspondence between subfields of M containing K and subgroups of Gal(M/K). Show
that if K ⊂ L ⊂ M then Gal(M/L) is a normal subgroup of Gal(M/K) if and only if
L/K is normal. What is Gal(L/K) in this case?
Let M be the splitting field of X 4 − 3 over Q. Prove that Gal(M/Q) is isomorphic
to the dihedral group of order 8. Hence determine all subfields of M , expressing each in
the form Q(x) for suitable x ∈ M .

Part II 2004
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B1/8 Differentiable Manifolds

What is a smooth vector bundle over a manifold M ?
Assuming the existence of “bump functions”, prove that every compact manifold
embeds in some Euclidean space Rn .
By choosing an inner product on Rn , or otherwise, deduce that for any compact
manifold M there exists some vector bundle η → M such that the direct sum T M ⊕ η is
isomorphic to a trivial vector bundle.

B2/7 Differentiable Manifolds

For each of the following assertions, either provide a proof or give and justify a
counterexample.
[You may use, without proof, your knowledge of the de Rham cohomology of
surfaces.]
(a) A smooth map f : S 2 → T 2 must have degree zero.
(b) An embedding ϕ : S 1 → Σg extends to an embedding ϕ̄ : D2 → Σg if and only if
the map Z
: H 1 (Σg ) → R
ϕ(S 1 )

is the zero map.

(c) RP1 × RP2 is orientable.
(d) The surface Σg admits the structure of a Lie group if and only if g = 1.

B4/4 Differentiable Manifolds

Define what it means for a manifold to be oriented, and define a volume form on
an oriented manifold.
Prove carefully that, for a closed connected oriented manifold of dimension n,
H n (M ) = R.
[You may assume the existence of volume forms on an oriented manifold.]
If M and N are closed, connected, oriented manifolds of the same dimension, define
the degree of a map f : M → N .
If f has degree d > 1 and y ∈ N , can f −1 (y) be
(i) infinite? (ii) a single point? (iii) empty?
Briefly justify your answers.

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B2/8 Algebraic Topology

Let K and L be finite simplicial complexes. Define the n-th chain group Cn (K)
and the boundary homomorphism dn : Cn (K) → Cn−1 (K). Prove that dn−1 dn = 0 and
define the homology groups of K. Explain briefly how a simplicial map f : K → L induces
a homomorphism f? of homology groups.
Suppose now that K consists of the proper faces of a 3-dimensional simplex.
Calculate from first principles the homology groups of K. If a simplicial map f : K → K
gives a homeomorphism of the underlying polyhedron |K|, is the induced homology map
f? necessarily the identity?

B3/7 Algebraic Topology

A finite simplicial complex K is the union of subcomplexes L and M . Describe
the Mayer-Vietoris exact sequence that relates the homology groups of K to those of L,
M and L ∩ M . Define all the homomorphisms in the sequence, proving that they are
well-defined (a proof of exactness is not required).
A surface X is constructed by identifying together (by means of a homeomorphism)
the boundaries of two Möbius strips Y and Z. Assuming relevant triangulations exist,
determine the homology groups of X.

B4/5 Algebraic Topology

Write down the definition of a covering space and a covering map. State and prove
the path lifting property for covering spaces and state, without proof, the homotopy lifting
property.
Suppose that a group G is a group of homeomorphisms of a space X. Prove that,
under conditions to be stated, the quotient map X → X/G is a covering map and that
π1 (X/G) is isomorphic to G. Give two examples in which this last result can be used to
determine the fundamental group of a space.

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B1/9 Number Fields

Let K = Q(θ), where θ is a root of X 3 − 4X + 1. Prove that K has degree 3 over
Q, and admits three distinct embeddings in R. Find the discriminant of K and determine
the ring of integers O of K. Factorise 2O and 3O into a product of prime ideals.
Using Minkowski’s bound, show that K has class number 1 provided all prime ideals
in O dividing 2 and 3 are principal. Hence prove that K has class number 1.
[You may assume that the discriminant of X 3 + aX + b is −4a3 − 27b2 .]

B2/9 Number Fields

Let m be an integer greater than 1 and let ζm denote a primitive m-th root of unity
in C. Let O be the ring of integers of Q(ζm ). If p is a prime number with (p, m) = 1,
outline the proof that
pO = ℘1 . . . ℘r ,
where the ℘i are distinct prime ideals of O, and r = ϕ(m)/f with f the least integer > 1
such that pf ≡ 1 mod m. [Here ϕ(m) is the Euler ϕ-function of m].
Determine the factorisations of 2, 3, 5 and 11 in Q(ζ5 ). For each integer n > 1,
prove that, in the ring of integers of Q(ζ5n ), there is a unique prime ideal dividing 2, and
a unique prime ideal dividing 3.

B4/6 Number Fields

Let K be a finite extension of Q, and O the ring of integers of K. Write an essay
outlining the proof that every non-zero ideal of O can be written as a product of non-zero
prime ideals, and that this factorisation is unique up to the order of the factors.

Part II 2004
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B1/10 Hilbert Spaces

Suppose that (en ) and (fm ) are orthonormal bases of a Hilbert space H and that
T ∈ L(H).
P∞ P∞ 2
(a) Show that n=1 kT (en )k2 = m=1 kT ∗ (fm )k .
P∞ P∞ 2
(b) Show that n=1 kT (en )k2 = m=1 kT (fm )k .
P∞
T ∈ L(H) is a Hilbert-Schmidt operator if n=1 kT (en )k2 < ∞ for some (and hence
every) orthonormal basis (en ).
(c) Show that the set
P∞HS of Hilbert-Schmidt operators forms a linear subspace of
L(H), and that hT, Si = n=1 hT (en ), S(en )i is an inner product on HS; show that this
inner product does not depend on the choice of the orthonormal basis (en ).
(d) Let kT kHS be the corresponding norm. Show that kT k 6 kT kHS , and show
that a Hilbert-Schmidt operator is compact.

B3/8 Hilbert Spaces

Let H be a Hilbert space. An operator T in L(H) is normal if T T ∗ = T ∗ T . Suppose
that T is normal and that σ(T ) ⊆ R. Let U = (T + iI)(T − iI)−1 .
(a) Suppose that A is invertible and AT = T A. Show that A−1 T = T A−1 .
(b) Show that U is normal, and that σ(U ) ⊆ {λ : |λ| = 1}.
(c) Show that U −1 is normal.
(d) Show that U is unitary.
(e) Show that T is Hermitian.
[You may use the fact that, if S is normal, the spectral radius of S is equal to kSk.]

B4/7 Hilbert Spaces

Suppose that T is a bounded linear operator on an infinite-dimensional Hilbert
space H, and that hT (x), xi is real and non-negative for each x ∈ H.
(a) Show that T is Hermitian.
(b) Let w(T ) = sup{hT (x), xi : kxk = 1}. Show that

kT (x)k2 6 w(T )hT (x), xi for each x ∈ H.

(c) Show that kT k is an approximate eigenvalue for T .

Suppose in addition that T is compact and injective.
(d) Show that kT k is an eigenvalue for T , with finite-dimensional eigenspace.
Explain how this result can be used to diagonalise T .

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B1/11 Riemann Surfaces

Let τ be a fixed complex number with positive imaginary part. For z ∈ C, define
∞
X
exp πiτ n2 + 2πin(z + 12 ) .


v(z) =
n=−∞

Prove the identities

v(z + 1) = v(z), v(−z) = v(z), v(z + τ ) = − exp(−πiτ − 2πiz) · v(z)

and deduce that v(τ /2) = 0. Show further that τ /2 is the only zero of v in the
parallelogram P with vertices −1/2, 1/2, 1/2 + τ , −1/2 + τ .
[You may assume that v is holomorphic on C.]

Now let {a1 , . . . , ak } and {b1 , . . . , bk } be two sets of complex numbers and

k
Y v(z − aj )
f (z) = .
j=1
v(z − bj )

Prove
Pk that f is a doubly-periodic meromorphic function, with periods 1 and τ , if and only
if j=1 (aj − bj ) is an integer.

B3/9 Riemann Surfaces

(a) Let f : R → S be a non-constant holomorphic map between compact connected
Riemann surfaces R and S.
Define the branching order vf (p) at a point p ∈ R and show that it is well-defined.
Show further that if h is a holomorphic map on S then vh◦f (p) = vh (f (p)) vf (p).
Define the degree of f and state the Riemann–Hurwitz formula. Show that if R has
Euler characteristic 0 then either S is the 2-sphere or vp (f ) = 1 for all p ∈ R.

(b) Let P and Q be complex polynomials of degree m ≥ 2 with no common roots.

Explain briefly how the rational function P (z)/Q(z) induces a holomorphic map F from
the 2-sphere S 2 ∼
= C ∪ {∞} to itself. What is the degree of F ? Show that there is at least
one and at most 2m − 2 points w ∈ S 2 such that the number of distinct solutions z ∈ S 2
of the equation F (z) = w is strictly less than deg F .

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B4/8 Riemann Surfaces

Let Λ be a lattice in C, Λ = Zω1 + Zω2 , where ω1 , ω2 6= 0 and ω1 /ω2 6∈ R. By
constructing an appropriate family of charts, show that the torus C/Λ is a Riemann surface
and that the natural projection π : z ∈ C → z + Λ ∈ C/Λ is a holomorphic map.
[You may assume without proof any known topological properties of C/Λ.]
Let Λ0 = Zω10 + Zω20 be another lattice in C, with ω10 , ω20 =
6 0 and ω10 /ω20 6∈ R.
By considering paths from 0 to an arbitrary z ∈ C, show that if f : C/Λ → C/Λ0 is a
conformal equivalence then

f (z + Λ) = (az + b) + Λ0 for some a, b, ∈ C, with a 6= 0.

[Any form of the Monodromy Theorem and basic results on the lifts of paths may be used
without proof, provided that these are correctly stated. You may assume without proof that
every injective holomorphic function F : C → C is of the form F (z) = az + b, for some
a, b ∈ C.]
Give an explicit example of a non-constant holomorphic map C/Λ → C/Λ that is
not a conformal equivalence.

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B2/10 Algebraic Curves

For each of the following curves C

(i) C = {(x, y) ∈ A2 |x3 − x = y 2 } (ii) C = {(x, y) ∈ A2 |x2 y + xy 2 = x4 + y 4 }

compute the points at infinity of C̄ ⊂ P2 (i.e. describe C̄ \ C), and find the singular points
of the projective curve C̄.
At which points of C̄ is the rational map C̄ 99K P1 , given by (X : Y : Z) 7→ (X : Y ),
not defined? Justify your answer.

B3/10 Algebraic Curves

(i) Let f : X → Y be a morphism of smooth projective curves. Define the divisor
f ∗ (D) if D is a divisor on Y , and state the “finiteness theorem”.
(ii) Suppose f : X → P1 is a morphism of degree 2, that X is smooth projective,
and that X 6= P1 . Let P, Q ∈ X be distinct ramification points for f . Show that, as
elements of cl(X), we have [P ] 6= [Q], but 2[P ] = 2[Q].

B4/9 Algebraic Curves

Let F (X, Y, Z) be an irreducible homogeneous polynomial of degree n, and write
C = {p ∈ P2 | F (p) = 0} for the curve it defines in P2 . Suppose C is smooth. Show that
the degree of its canonical class is n(n − 3).
Hence, or otherwise, show that a smooth curve of genus 2 does not embed in P2 .

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B1/13 Probability and Measure

Let (Ω, F, P) be a probability space and let A = (Ai : i = 1, 2, . . .) be a sequence
of events.
(a) What is meant by saying that A is a family of independent events?
(b) Define the events {An infinitely often} and {An eventually}. State and prove
the two Borel–Cantelli lemmas for A.
(c) Let X1 , X2 , . . . be the outcomes of a sequence of independent flips of a fair coin,

1
P(Xi = 0) = P(Xi = 1) = 2 for i > 1.

Let Ln be the length of the run beginning at the nth flip. For example, if the first fourteen
outcomes are 01110010000110, then L1 = 1, L2 = 3, L3 = 2, etc.
Show that  
Ln
P lim sup >1 =0,
n→∞ log2 n

and furthermore that  

Ln
P lim sup =1 =1.
n→∞ log2 n

B2/12 Probability and Measure

Let (Ω, F, µ) be a measure space and let 1 6 p 6 ∞.
(a) Define the Lp -norm ||f ||p of a measurable function f : Ω → R, and define the space
Lp (Ω, F, µ).
(b) Prove Minkowski’s inequality:

||f + g||p 6 ||f ||p + ||g||p for f, g ∈ Lp (Ω, F, µ), 1 6 p 6 ∞.

[You may use Hölder’s inequality without proof provided it is clearly stated.]
(c) Explain what is meant by saying that Lp (Ω, F, µ) is complete. Show that
∞
L (Ω, F, µ) is complete.
(d) Suppose that {fn : n > 1} is a sequence of measurable functions satisfying
||fn ||p → 0 as n → ∞.
(i) Show that if p = ∞, then fn → 0 almost everywhere.
(ii) When 1 6 p < ∞, give an example of a measure space (Ω, F, µ) and such a
sequence {fn } such that, for all ω ∈ Ω, fn (ω) 6→ 0 as n → ∞.

Part II 2004
 77

B3/12 Probability and Measure

(a) Let (Ω, F, P) be a probability space and let θ : Ω → Ω be measurable. What is
meant by saying that θ is measure-preserving? Define an invariant event and an
invariant random variable, and explain what is meant by saying that θ is ergodic.
(b) Let m be a probability measure on (R, B). Let

Ω = RN = {x = (x1 , x2 , . . .) : xi ∈ R for i > 1} ,

let F be the smallest σ-field of Ω with respect to which the coordinate maps
Xn (x) = xn , for x ∈ Ω, n > 1, are measurable, and let P be the unique probability
measure on (Ω, F) satisfying
n
Y
P(Xi ∈ Ai for 1 6 i 6 n) = m(Ai )
i=1

for all Ai ∈ B, n > 1. Define θ : Ω → Ω by θ(x) = (x2 , x3 , . . .) for x = (x1 , x2 , . . .).

(i) Show that θ is measurable and measure-preserving.
(ii) Define the tail σ-field T of the coordinate maps X1 , X2 , . . ., and show that
the invariant σ-field I of θ satisfies I ⊆ T . Deduce that θ is ergodic. [Any
general result used must be stated clearly but the proof may be omitted.]
(c) State Birkhoff’s ergodic theorem and explain how to deduce that, given independent
identically-distributed integrable random variables Y1 , Y2 , . . ., there exists ν ∈ R
such that
1
(Y1 + Y2 + · · · + Yn ) → ν a.e. as n → ∞ .
n

B4/11 Probability and Measure

Let (Ω, F, P) be a probability space and let X, X1 , X2 , . . . be random variables.
Write an essay in which you discuss the statement: if Xn → X almost everywhere, then
E(Xn ) → E(X). You should include accounts of monotone, dominated, and bounded
convergence, and of Fatou’s lemma.
[You may assume without proof the following fact. Let (Ω, F, µ) be a measure space,
and let f : Ω → R be non-negative with finite integral µ(f ). If (fn : n > 1) are non-negative
measurable functions with fn (ω) ↑ f (ω) for all ω ∈ Ω, then µ(fn ) → µ(f ) as n → ∞.]

Part II 2004
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B2/13 Applied Probability

Let M be a Poisson random measure of intensity λ on the plane R2 . Denote by
C(r) the circle {x ∈ R2 : ||x|| < r} of radius r in R2 centred at the origin and let Rk be
the largest radius such that C(Rk ) contains precisely k points of M . [Thus C(R0 ) is the
largest circle about the origin containing no points of M , C(R1 ) is the largest circle about
the origin containing a single point of M , and so on.] Calculate ER0 , ER1 and ER2 .
Now let N be a Poisson random measure of intensity λ on the line R1 . Let Lk
be the length of the largest open interval that covers the origin and contains precisely k
points of N . [Thus L0 gives the length of the largest interval containing the origin but no
points of N, L1 gives the length of the largest interval containing the origin and a single
point of N , and so on.] Calculate EL0 , EL1 and EL2 .

B3/13 Applied Probability

Let (Xt , t ≥ 0) be a renewal process with holding times (Sn , n = 1, 2, . . .) and
(Yt , t ≥ 0) be a renewal-reward process over (Xt ) with a sequence of rewards
(Wn , n = 1, 2, . . .). Under assumptions on (Sn ) and (Wn ) which you should state clearly,
prove that the ratios
Xt /t and Yt /t
converge as t → ∞. You should specify the form of convergence guaranteed by your
assumptions. The law of large numbers, in the appropriate form, for sums S1 + . . . + Sn
and W1 + . . . + Wn can be used without proof.
In a mountain resort, when you rent skiing equipment you are given two options.
(1) You buy an insurance waiver that costs C/4 where C is the daily equipment rent.
Under this option, the shop will immediately replace, at no cost to you, any piece of
equipment you break during the day, no matter how many breaks you had. (2) If you
don’t buy the waiver, you’ll pay 3C in the case of any break.
To find out which option is better for me, I decided to set up two models of renewal-
reward process (Yt ). In the first model, (Option 1), all of the holding times Sn are equal
to 6. In the second model, given that there is no break on day n (an event of probability
4/5), we have Sn = 6, Wn = C, but given that there is a break on day n, we have that
Sn is uniformly distributed on (0, 6), and Wn = 4C. (In the second model, I would not
continue skiing after a break, whereas in the first I would.)
Calculate in each of these models the limit

lim Yt /t
t→∞

representing the long-term average cost of a unit of my skiing time.

Part II 2004
 79

B4/12 Applied Probability

Consider an M/G/1 queue with ρ = λES < 1. Here λ is the arrival rate and ES is
the mean service time. Prove that in equilibrium, the customer’s waiting time W has the
moment-generating function MW (t) = E etW given by

(1 − ρ)t
MW (t) =
t + λ(1 − MS (t))

where MS (t) = EetS is the moment-generating function of service time S.

[You may assume that in equilibrium, the M/G/1 queue size X at the time
immediately after the customer’s departure has the probability generating function

(1 − ρ)(1 − z)MS (λ(z − 1))

E zX = , 0 6 z < 1 .]
MS (λ(z − 1)) − z

Deduce that when the service times are exponential of rate µ then

λ(1 − ρ)
MW (t) = 1 − ρ + , −∞ < t < µ − λ .
µ−λ−t

Further, deduce that W takes value 0 with probability 1 − ρ and that

P(W > x|W > 0) = e−(µ−λ)x , x > 0.

Sketch the graph of P(W > x) as a function of x.

Now consider the M/G/1 queue in the heavy traffic approximation, when the
service-time distribution is kept fixed and the arrival rate λ → 1/ES, so that ρ → 1.
Assuming that the second moment ES 2 < ∞, check that the limiting distribution of the
re-scaled waiting time W̃λ = (1 − λES)W is exponential, with rate 2ES/ES 2 .

Part II 2004
 80

B1/14 Information Theory

State the formula for the capacity of a memoryless channel.
(a) Consider a memoryless channel where there are two input symbols, A and B, and
three output symbols, A, B, ∗. Suppose each input symbol is left intact with probability
1/2, and transformed into a ∗ with probability 1/2. Write down the channel matrix, and
calculate the capacity.
(b) Now suppose the output is further processed by someone who cannot distinguish
A and ∗, so that the matrix becomes
 
1 0
.
1/2 1/2

Calculate the new capacity.

B2/14 Information Theory

For integer-valued random variables X and Y , define the relative entropy hY (X)
of X relative to Y .
Prove that hY (X) > 0, with equality if and only if P(X = x) = P(Y = x) for all x.
By considering Y , a geometric random variable with parameter chosen appropri-
ately, show that if the mean EX = µ < ∞, then
h(X) 6 (µ + 1) log(µ + 1) − µ log µ ,
with equality if X is geometric.

B4/13 Information Theory

Define a cyclic code of length N .
Show how codewords can be identified with polynomials in such a way that cyclic
codes correspond to ideals in the polynomial ring with a suitably chosen multiplication
rule.
Prove that any cyclic code X has a unique generator, i.e. a polynomial c(X) of
minimum degree, such that the code consists of the multiples of this polynomial. Prove
that the rank of the code equals N − deg c(X), and show that c(X) divides X N + 1.
Let X be a cyclic code. Set
N
X
X ⊥ = {y = y1 . . . yN : xi yi = 0 for all x = x1 . . . xN ∈ X }
i=1

(the dual code). Prove that X ⊥ is cyclic and establish how the generators of X and X ⊥
are related to each other.
Show that the repetition and parity codes are cyclic, and determine their genera-
tors.

Part II 2004
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B2/15 Optimization and Control

A gambler is presented with a sequence of n > 6 random numbers, N1 , N2 , . . . , Nn ,
one at a time. The distribution of Nk is

P (Nk = k) = 1 − P (Nk = −k) = p ,

where 1/(n − 2) < p ≤ 1/3. The gambler must choose exactly one of the numbers, just
after it has been presented and before any further numbers are presented, but must wait
until all the numbers are presented before his payback can be decided. It costs £1 to play
the game. The gambler receives payback as follows: nothing if he chooses the smallest of
all the numbers, £2 if he chooses the largest of all the numbers, and £1 otherwise.
Show that there is an optimal strategy of the form “Choose the first number k such
that either (i) Nk > 0 and k ≥ n − r0 , or (ii) k = n − 1”, where you should determine the
constant r0 as explicitly as you can.

B3/14 Optimization and Control

The strength of the economy evolves according to the equation

ẍt = −α2 xt + ut ,

where x0 = ẋ0 = 0 and ut is the effort that the government puts into reform at time
t, t ≥ 0. The government wishes to maximize its chance of re-election at a given future
time T , where this chance is some monotone increasing function of
Z T
1
xT − u2t dt .
2 0

Use Pontryagin’s maximum principle to determine the government’s optimal reform

policy, and show that the optimal trajectory of xt is

t −2 1
xt = α cos(α(T − t)) − α−3 cos(αT ) sin(αt) .
2 2

Part II 2004
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B4/14 Optimization and Control

Consider the deterministic dynamical system

ẋt = Axt + But

where A and B are constant matrices, xt ∈ Rn , and ut is the control variable, ut ∈ Rm .

What does it mean to say that the system is controllable?
Let yt = e−tA xt − x0 . Show that if Vt is the set of possible values for yt as the
control {us : 0 ≤ x ≤ t} is allowed to vary, then Vt is a vector space.
Show that each of the following three conditions is equivalent to controllability of
the system.
(i) The set {v ∈ Rn : v > yt = 0 for all yt ∈ Vt } = {0}.
Rt >
(ii) The matrix H(t) = 0 e−sA BB > e−sA ds is (strictly) positive definite.
(iii) The matrix Mn = [B AB A2 B · · · An−1 B] has rank n.
Consider the scalar system
n
X  d n−j
aj ξt = ut ,
j=0
dt

where a0 = 1. Show that this system is controllable.

Part II 2004
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B1/18 Partial Differential Equations

(a) State and prove the Mean Value Theorem for harmonic functions.
(b) Let u > 0 be a harmonic function on an open set Ω ⊂ Rn . Let B(x, a) = {y ∈
n
R : |x − y| < a}. For any x ∈ Ω and for any r > 0 such that B(x, 4r) ⊂ Ω, show that

sup u(y) 6 3n inf u(y) .

{y∈B(x,r)} {y∈B(x,r)}

B2/18 Partial Differential Equations

(a) State and prove the Duhamel principle for the wave equation.
(b) Let u ∈ C 2 ([0, T ] × Rn ) be a solution of

utt + ut − ∆u + u = 0

where ∆ is taken in the variables x ∈ Rn and ut = ∂t u etc.

Using an ‘energy method’, or otherwise, show that, if u = ut = 0 on the set
{t = 0, |x − x0 | 6 t0 } for some (t0 , x0 ) ∈ [0, T ] × Rn , then u vanishes on the region
K(t, x) = {(t, x) : 0 6 t 6 t0 , |x − x0 | 6 t0 − t}. Hence deduce uniqueness for the Cauchy
problem for the above PDE with Schwartz initial data.

Part II 2004
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B3/18 Partial Differential Equations

(i) Find w : [0, ∞) × R −→ R such that w(t, ·) is a Schwartz function of ξ for each
t and solves
wt (t, ξ) + (1 + ξ 2 )w(t, ξ) = g(ξ) ,
w(0, ξ) = w0 (ξ) ,
where g and w0 are given Schwartz functions and wt denotes ∂t w. If F represents the
Fourier transform operator in the ξ variables only and F −1 represents its inverse, show
that the solution w satisfies

∂t (F −1 )w(t, x) = F −1 (∂t w)(t, x)

and calculate lim w(t, ·) in Schwartz space.

t→∞

(ii) Using the results of Part (i), or otherwise, show that there exists a solution of
the initial value problem

ut (t, x) − uxx (t, x) + u(t, x) = f (x)

u(0, x) = u0 ,
with f and u0 given Schwartz functions, such that

ku(t, ·) − φkL∞ (R) −→ 0

as t → ∞ in Schwartz space, where φ is the solution of

00
−φ + φ = f.

B4/18 Partial Differential Equations

(a) State a theorem of local existence, uniqueness and C 1 dependence on the initial
data for a solution for an ordinary differential equation. Assuming existence, prove that
the solution depends continuously on the initial data.
(b) State a theorem of local existence of a solution for a general quasilinear first–
order partial differential equation with data on a smooth non-characteristic hypersurface.
Prove this theorem in the linear case assuming the validity of the theorem in part (a);
explain in your proof the importance of the non-characteristic condition.

Part II 2004
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B1/19 Methods of Mathematical Physics

State the convolution theorem for Laplace transforms.
The temperature T (x, t) in a semi-infinite rod satisfies the heat equation

∂2T 1 ∂T
= , x ≥ 0, t ≥ 0
∂x2 k ∂t

and the conditions T (x, 0) = 0 for x ≥ 0, T (0, t) = f (t) for t ≥ 0 and T (x, t) → 0 as
x → ∞. Show that Z t
T (x, t) = f (τ ) G(x, t − τ )dτ,
0

where r
x2 2
G(x, t) = 3
e−x /4kt .
4πkt

Part II 2004
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B2/19 Methods of Mathematical Physics

(a) The Beta function is defined by
Z 1
B (p, q) = xp−1 (1 − x)q−1 dx .
0

Show that Z ∞
B (p, q) = x−p−q (x − 1)q−1 dx .
1

(b) The function J(p, q) is defined by

Z
J(p, q) = tp−1 (1 − t)q−1 dt ,
γ

where the integrand has a branch cut along the positive real axis. Just above the cut,
arg t = 0. For t > 1 just above the cut, arg (1 − t) = −π. The contour γ runs from
t = ∞e2πi , round the origin in the negative sense, to t = ∞ (i.e. the contour is a reflection
of the usual Hankel contour). What restriction must be placed on p and q for the integral
to converge?
By evaluating J(p, q) in two ways, show that

1 − e2πip B (p, q) + e−πi(q−1) − eπi(2p+q−1) B (1 − p − q, q) = 0 ,

where p and q are any non-integer complex numbers.

Using the identity
Γ(p)Γ(q)
B(p, q) = ,
Γ(p + q)
deduce that

Γ(p)Γ(1 − p) sin(πp) = Γ(p + q)Γ(1 − p − q) sin[π(1 − p − q)] ,

and hence that

π = Γ(q)Γ(1 − q) sin[π(1 − q)] .

Part II 2004
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B3/19 Methods of Mathematical Physics

The function w(z) satisfies the third-order differential equation

d3 w
− zw = 0
dz 3

subject to the conditions w(z) → 0 as z → ±i∞ and w(0) = 1 . Obtain an integral

representation for w(z) of the form
Z
w(z) = ezt f (t)dt ,
γ

and determine the function f (t) and the contour γ.

Using the change of variable t = z 1/3 τ , or otherwise, compute the leading term in
the asymptotic expansion of w(z) as z → +∞.

B4/19 Methods of Mathematical Physics

Let h(t) = i(t + t2 ) . Sketch the path of Im(h(t)) = const. through the point t = 0,
and the path of Im(h(t)) = const. through the point t = 1 .
By integrating along these paths, show that as λ → ∞
1
c2 e2iλ
Z
2 c1
t−1/2 eiλ(t+t ) dt ∼ + ,
0 λ1/2 λ

where the constants c1 and c2 are to be computed.

Part II 2004
 88

B1/21 Electrodynamics
The Maxwell field tensor is

0 −Ex −Ey −Ez

 
 Ex 0 −Bz By 
F ab = ,
Ey Bz 0 −Bx
Ez −By Bx 0

and the 4-current density is J a = (ρ, j). Write down the 3-vector form of Maxwell’s
equations and the continuity equation, and obtain the equivalent 4-vector equations.
Consider a Lorentz transformation from a frame F to a frame F 0 moving with
relative (coordinate) velocity v in the x-direction

γ γv 0 0
 
 γv γ 0 0
La b = ,
0 0 1 0
0 0 0 1
√
where γ = 1/ 1 − v 2 . Obtain the transformation laws for E and B. Which quantities,
quadratic in E and B, are Lorentz scalars?

B2/21 Electrodynamics
A particle of rest mass m and charge q moves along a path xa (s), where s is the
particle’s proper time. The equation of motion is

mẍa = qF ab ηbc ẋc ,

where ẋa = dxa /ds etc., F ab is the Maxwell field tensor (F 01 = −Ex , F 23 = −Bx ,
where Ex and Bx are the x-components of the electric and magnetic fields) and ηbc is the
Minkowski metric tensor. Show that ẋa ẍa = 0 and interpret both the equation of motion
and this equation in the classical limit.
The electromagnetic field is given in cartesian coordinates by E = (0, E, 0) and
B = (0, 0, E), where E is constant and uniform. The particle starts from rest at the
origin. Show that the orbit is given by

9x2 = 2αy 3 , z = 0,

where α = qE/m.

Part II 2004
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B4/21 Electrodynamics
Using Lorentz gauge, Aa ,a = 0, Maxwell’s equations for a current distribution J a
can be reduced to Aa (x) = µ0 J a (x). The retarded solution is
Z
µ0
Aa (x) = d4 y θ(z 0 )δ(zc z c )J a (y),
2π

where z a = xa − y a . Explain, heuristically, the rôle of the δ-function and Heaviside step
function θ in this formula.
The current distribution is produced by a point particle of charge q moving on a
world line ra (s), where s is the particle’s proper time, so that
Z
a
J (y) = q ds V a (s)δ (4) (y − r(s)),

where V a = ṙa (s) = dra /ds. Show that

Z
µ0 q
Aa (x) = ds θ(X 0 )δ(Xc X c )V a (s),
2π

where X a = xa − ra (s), and further that, setting α = Xc V c ,

µ0 q V a
 
a
A (x) = ,
4π α s=s∗

where s∗ should be defined. Verify that

 
∗ Xa
s ,a = .
α s=s∗

Evaluating quantities at s = s∗ show that

 a
V 1
= 2 [−V a Vb + S a Xb ] ,
α ,b α

where S a = V̇ a + V a (1 − Xc V̇ c )/α. Hence verify that Aa ,a (x) = 0 and

µ0 q
Fab = (Sa Xb − Sb Xa ) .
4πα2

Verify this formula for a stationary point charge at the origin.

[Hint: If f (s) has simple zeros at si , i = 1, 2, . . . then

X δ(si )
δ(f (s)) = .
i
|f 0 (si )|

Part II 2004
 90

B1/22 Statistical Physics

Define the notions of entropy S and thermodynamic temperature T for a gas of
particles in a variable volume V . Derive the fundamental relation

dE = T dS − P dV .

The free energy of the gas is defined as F = E − T S. Why is it convenient to regard

F as a function of T and V ? By considering F , or otherwise, show that

∂S ∂P
= .
∂V ∂T
T V

Deduce that the entropy of an ideal gas, whose equation of state is P V = N T (using
energy units), has the form
 
V
S = N log + N c(T ) ,
N

where c(T ) is independent of N and V .

Show that if the gas is in contact with a heat bath at temperature T , then the
probability of finding the gas in a particular quantum microstate of energy Er is

Pr = e(F −Er )/T .

Part II 2004
 91

B3/22 Statistical Physics

Describe briefly why a low density gas can be investigated using classical statistical
mechanics.
Explain why, for a gas of N structureless atoms, the measure on phase space is

1 d3N q d3N p
,
N ! (2π~)3N

and the probability density in phase space is proportional to

e−E(q,p)/T ,

where T is the temperature in energy units.

Derive the Maxwell probability distribution for atomic speeds v,
 m 3/2 2
ρ(v) = 4πv 2 e−mv /2T .
2πT
Why is this valid even if the atoms interact?
Find the mean value v̄ of the speed of the atoms.
Is 21 m(v̄)2 the mean kinetic energy of the atoms?

B4/23 Statistical Physics

Derive the Bose-Einstein expression for the mean number of Bose particles n̄
occupying a particular single-particle quantum state of energy ε, when the chemical
potential is µ and the temperature is T in energy units.
Why is the chemical potential for a gas of photons given by µ = 0?
Show that, for black-body radiation in a cavity of volume V at temperature T , the
mean number of photons in the angular frequency range (ω, ω + dω) is

V ω 2 dω
.
π 2 c3 e~ω/T − 1

Hence, show that the total energy E of the radiation in the cavity is

E = KV T 4 ,

where K is a constant that need not be evaluated.

Use thermodynamic reasoning to find the entropy S and pressure P of the radiation
and verify that
E − TS + PV = 0 .
Why is this last result to be expected for a gas of photons?

Part II 2004
 92

B1/23 Applications of Quantum Mechanics

The operator corresponding to a rotation through an angle θ about an axis n, where
n is a unit vector, is
U (n, θ) = eiθ n·J/~ .
If U is unitary show that J must be hermitian. Let V = (V1 , V2 , V3 ) be a vector operator
such that
U (n, δθ)VU (n, δθ)−1 = V + δθ n × V .
Work out the commutators [Ji , Vj ]. Calculate

U (ẑ, θ)VU (ẑ, θ)−1 ,

for each component of V.

If |jmi are standard angular momentum states determine hjm0 |U (ẑ, θ)|jmi for any
j, m, m0 and also determine h 21 m0 |U (ŷ, θ)| 12 mi.

Hint : J3 |jmi = m~|jmi, J+ | 21 − 12 i = ~| 12 12 i, J− | 12 21 i = ~| 12 − 21 i.

 

Part II 2004
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B2/23 Applications of Quantum Mechanics

The wave function for a single particle with a potential V (r) has the asymptotic
form for large r
eikr
ψ(r, θ) ∼ eikr cos θ + f (θ) .
r
How is f (θ) related to observable quantities? Show how f (θ) can be expressed in terms
of phase shifts δ` (k) for ` = 0, 1, 2, . . ..
Assume that V (r) = 0 for r ≥ a, and let R` (r) denote the solution of the radial
Schrödinger equation, regular at r = 0, with energy ~2 k 2 /2m and angular momentum `.
Let N` (k) = aR` 0 (a)/R` (a). Show that

N` (k) j` (ka) − ka j` 0 (ka)

tan δ` (k) = .
N` (k) n` (ka) − ka n` 0 (ka)

Assuming that N` (k) is a smooth function for k ≈ 0, determine the expected behaviour
of δ` (k) as k → 0. Show that for k → 0 then f (θ) → c, with c a constant, and determine
c in terms of N0 (0).


For V = 0 the two independent solutions of the radial Schrödinger equation are j` (kr)
and n` (kr) with

1 1
j` (ρ) ∼ sin(ρ − 12 `π), n` (ρ) ∼ − cos(ρ − 12 `π) as ρ → ∞ ,
ρ ρ
j` (ρ) ∝ ρ` , n` (ρ) ∝ ρ−`−1 as ρ → 0 ,
X∞
iρ cos θ
e = (2` + 1)i` j` (ρ) P` (cos θ) ,
`=0
sin ρ cos ρ
j0 (ρ) = , n0 (ρ) = − .
ρ ρ

Part II 2004
 94

B3/23 Applications of Quantum Mechanics

For a periodic potential V (r) = V (r + `), where ` is a lattice vector, show that we
may write X
V (r) = ag eig·r , ag ∗ = a−g ,
g

where the set of g should be defined.

Show how to construct general wave functions satisfying ψ(r + `) = eik·` ψ(r) in terms of
free plane-wave wave-functions.
Show that the nearly free electron model gives an energy gap 2|ag | when k = 12 g.
Explain why, for a periodic potential, the allowed energies form bands En (k) where k may
be restricted to a single Brillouin zone. Show that En (k) = En (k + g) if k and k + g
belong to the Brillouin zone.
How are bands related to whether a material is a conductor or an insulator?

Part II 2004
 95

B4/24 Applications of Quantum Mechanics

Describe briefly the variational approach to determining approximate energy eigen-
values for a Hamiltonian H.
Consider a Hamiltonian H and two states |ψ1 i, |ψ2 i such that

hψ1 |H|ψ1 i = hψ2 |H|ψ2 i = E , hψ2 |H|ψ1 i = hψ1 |H|ψ2 i = ε ,

hψ1 |ψ1 i = hψ2 |ψ2 i = 1 , hψ2 |ψ1 i = hψ1 |ψ2 i = s .

Show that, by considering a linear combination α |ψ1 i + β |ψ2 i, the variational method
gives
E −ε E +ε
, ,
1−s 1+s
as approximate energy eigenvalues.
Consider the Hamiltonian for an electron in the presence of two protons at 0 and R,

p2 e2
 
1 1 1
H= + − − , R = |R| .
2m 4π0 R |r| |r − R|
1
Let ψ0 (r) = e−r/a /(πa3 ) 2 be the ground state hydrogen atom wave function which satisfies
 p2 e2 
− ψ0 (r) = E0 ψ0 (r) .
2m 4π0 |r|

It is given that

R2
Z 
R
S = d r ψ0 (r)ψ0 (r − R) = 1 + + 2 e−R/a ,
3
a 3a
Z  
1 1 R −R/a
U = d3 r ψ0 (r)ψ0 (r − R) = 1+ e ,
|r| a a

and, for large R, that

Z
1 1
d3 r ψ0 (r)2 − = O e−2R/a .


|r − R| R

Consider the trial wave function α ψ0 (r) + β ψ0 (r − R). Show that the variational estimate
for the ground state energy for large R is

e2
S − RU ) + O e−2R/a .


E(R) = E0 +
4π0 R

Explain why there is an attractive force between the two protons for large R.

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B1/25 Fluid Dynamics II

Consider a uniform stream of inviscid incompressible fluid incident onto a two-
dimensional body (such as a circular cylinder). Sketch the flow in the region close to the
stagnation point, S, at the front of the body.
Let the fluid now have a small but non-zero viscosity. Using local co-ordinates x
along the boundary and y normal to it, with the stagnation point as origin and y > 0 in
the fluid, explain why the local outer, inviscid flow is approximately of the form

u = (Ex, −Ey)

for some positive constant E.

Use scaling arguments to find the thickness δ of the boundary layer on the body
near S. Hence show that there is a solution of the boundary layer equations of the form

u(x, y) = Exf 0 (η),

where η is a suitable similarity variable and f satisfies

2
f 000 + f f 00 − f 0 = −1. (∗)

What are the appropriate boundary conditions for (∗) and why? Explain briefly how you
would obtain a numerical solution to (∗) subject to the appropriate boundary conditions.
Explain why it is neither possible nor appropriate to perform a similar analysis
near the rear stagnation point of the inviscid flow.

B2/25 Fluid Dynamics II

An incompressible fluid with density ρ and viscosity µ is forced by a pressure
difference ∆p through the narrow gap between two parallel circular cylinders of radius a
with axes 2a + b apart. Explaining any approximations made, show that, provided b  a
and ρb3 ∆p  µ2 a, the volume flux (per unit length of cylinder) is

2b5/2 ∆p
9πa1/2 µ

when the cylinders are stationary.

Show also that when the two cylinders rotate with angular velocities Ω and −Ω
respectively, the change in the volume flux is

4
baΩ.
3

For the case ∆p = 0, find and sketch the function f (x) = u0 (x)/(aΩ), where u0 is the
centreline velocity at position x along the gap in the direction of flow. Comment on the
values taken by f .

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B3/24 Fluid Dynamics II

Using the Milne-Thompson circle theorem, or otherwise, write down the complex
potential w describing inviscid incompressible two-dimensional flow past a circular cylinder
of radius a centred on the origin, with circulation κ and uniform velocity (U, V ) in the far
field.
Hence, or otherwise, find an expression for the velocity field if the cylinder is
replaced by a flat plate of length 4a, centred on the origin and aligned with the x-axis.
Evaluate the velocity field on the two sides of the plate and confirm that the normal
velocity is zero.
Explain the significance of the Kutta condition, and determine the value of the
circulation that satisfies the Kutta condition when U > 0.
With this value of the circulation, calculate the difference in pressure between the
upper and lower sides of the plate at position x (−2a ≤ x ≤ 2a). Comment briefly on the
value of the pressure at the leading edge and the force that this would produce if the plate
had a small non-zero thickness.
Determine the force on the plate, explaining carefully the direction in which it acts.
I  2
iρ dw
[The Blasius formula Fx − iFy = dz, where C is a closed contour lying just
2 C dz
outside the body, may be used without proof.]

B4/26 Fluid Dynamics II

Write an essay on the Kelvin-Helmholtz instability of a vortex sheet. Your essay
should include a detailed linearised analysis, a physical interpretation of the instability,
and an informal discussion of nonlinear effects and of the effects of viscosity.

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B1/26 Waves in Fluid and Solid Media

A physical system permits one-dimensional wave propagation in the x-direction
according to the equation
∂2ψ ∂6ψ
2
− α2 6 = 0 ,
∂t ∂x
where α is a real positive constant. Derive the corresponding dispersion relation and sketch
graphs of frequency, phase velocity and group velocity as functions of the wave number.
Is it the shortest or the longest waves that are at the front of a dispersing wave train
arising from a localised initial disturbance? Do the wave crests move faster or slower than
a packet of waves?
Find the solution of the above equation for the initial disturbance given by
Z ∞
∂ψ
ψ(x, 0) = A(k)eikx dk , (x, 0) = 0 ,
−∞ ∂t

where A(k) is real and A(−k) = A(k).

Use the method of stationary phase to obtain a leading-order approximation to this
solution for large t when V = x/t is held fixed.
[Note that Z ∞
2 1
e±iu du = π 2 e±iπ/4 . ]
−∞

B2/26 Waves in Fluid and Solid Media

The linearised equation of motion governing small disturbances in a homogeneous
elastic medium of density ρ is

∂2u
ρ = (λ + µ)∇(∇ · u) + µ∇2 u ,
∂t2

where u(x, t) is the displacement, and λ and µ are the Lamé constants. Derive solutions for
plane longitudinal waves P with wavespeed cP , and plane shear waves S with wavespeed
cS .
The half-space y < 0 is filled with the elastic solid described above, while the slab
0 < y < h is filled with an elastic solid with shear modulus µ, and wavespeeds cP and
cS . There is a vacuum in y > h. A harmonic plane SH wave of frequency ω and unit
amplitude propagates from y < 0 towards the interface y = 0. The wavevector is in the
xy-plane, and makes an angle θ with the y-axis. Derive the complex amplitude, R, of the
reflected SH wave in y < 0. Evaluate |R| for all possible values of cS /cS , and explain your
answer.

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B3/25 Waves in Fluid and Solid Media

The dispersion relation for sound waves of frequency ω in a stationary, homogeneous
gas is ω = c|k|, where c is the speed of sound and k is the wavevector. Derive the dispersion
relation for sound waves of frequency ω in a uniform flow with velocity U.
For a slowly-varying medium with a local dispersion relation ω = Ω(k; x, t), derive
the ray-tracing equations

dxi ∂Ω dki ∂Ω dω ∂Ω
= , =− , = .
dt ∂ki dt ∂xi dt ∂t

The meaning of the notation d/dt should be carefully explained.

Suppose that two-dimensional sound waves with initial wavevector (k0 , l0 ) are
generated at the origin in a gas occupying the half-space y > 0. The gas has a mean
1
velocity (γy, 0), where 0 < γ  (k02 + l02 ) 2 . Show that
(a) if k0 > 0 and l0 > 0 then an initially upward propagating wavepacket returns to
the level y = 0 within a finite time, after having reached a maximum height that
should be identified;
(b) if k0 < 0 and l0 > 0 then an initially upward propagating wavepacket continues to
propagate upwards for all time.
For the case of a fixed frequency disturbance comment briefly on whether or not there is
a quiet zone.

B4/27 Waves in Fluid and Solid Media

A plane shock is moving with speed U into a perfect gas. Ahead of the shock the gas
is at rest with pressure p1 and density ρ1 , while behind the shock the velocity, pressure and
density of the gas are u2 , p2 and ρ2 respectively. Derive the Rankine-Hugoniot relations
across the shock. Show that
ρ1 2c2 + (γ − 1)U 2
= 1 ,
ρ2 (γ + 1)U 2
where c21 = γp1 /ρ1 and γ is the ratio of the specific heats of the gas. Now consider a
change of frame such that the shock is stationary and the gas has a component of velocity
V parallel to the shock. Deduce that the angle of deflection δ of the flow which is produced
by a stationary shock inclined at an angle α = tan−1 (U/V ) to an oncoming stream of Mach
1
number M = (U 2 + V 2 ) 2 /c1 is given by

2 cot α(M 2 sin α2 − 1)

tan δ = .
2 + M 2 (γ + cos 2α)

[Note that
tan θ + tan φ
tan(θ + φ) = . ]
1 − tan θ tan φ

Part II 2004