MATHEMATICAL TRIPOS Part II 2018

List of Courses

Algebraic Geometry
Algebraic Topology
Analysis of Functions
Applications of Quantum Mechanics
Applied Probability
Asymptotic Methods
Automata and Formal Languages
Classical Dynamics
Coding & Cryptography
Cosmology
Differential Geometry
Dynamical Systems
Electrodynamics
Fluid Dynamics II
Further Complex Methods
Galois Theory
General Relativity
Graph Theory
Integrable Systems
Linear Analysis
Logic and Set Theory
Mathematical Biology
Number Fields
Number Theory
Numerical Analysis
Optimisation and Control
Principles of Quantum Mechanics
Principles of Statistics
Probability and Measure
Quantum Information and Computation

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Representation Theory
Riemann Surfaces
Statistical Modelling
Statistical Physics
Stochastic Financial Models
Topics in Analysis
Waves

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Paper 4, Section II
24I Algebraic Geometry
State a theorem which describes the canonical divisor of a smooth plane curve C in
terms of the divisor of a hyperplane section. Express the degree of the canonical divisor
KC and the genus of C in terms of the degree of C. [You need not prove these statements.]
From now on, we work over C. Consider the curve in A2 defined by the equation

y + x3 + xy 3 = 0.

Let C be its projective completion. Show that C is smooth.

Compute the genus of C by applying the Riemann–Hurwitz theorem to the mor-
phism C → P1 induced from the rational map (x, y) 7→ y. [You may assume that the
discriminant of x3 + ax + b is −4a3 − 27b2 .]

Paper 3, Section II
24I Algebraic Geometry
(a) State the Riemann–Roch theorem.
(b) Let E be a smooth projective curve of genus 1 over an algebraically closed field
k, with char k 6= 2, 3. Show that there exists an isomorphism from E to the plane cubic
in P2 defined by the equation

y 2 = (x − λ1 )(x − λ2 )(x − λ3 ),

for some distinct λ1 , λ2 , λ3 ∈ k.

(c) Let Q be the point at infinity on E. Show that the map E → Cl0 (E), P 7→ [P −Q]
is an isomorphism.
Describe how this defines a group structure on E. Denote addition by ⊞. Determine
all the points P ∈ E with P ⊞P = Q in terms of the equation of the plane curve in part (b).

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Paper 2, Section II
24I Algebraic Geometry
(a) Let X ⊆ An be an affine algebraic variety defined over the field k.
Define the tangent space Tp X for p ∈ X, and the dimension of X in terms of Tp X.
Suppose that k is an algebraically closed field with char k > 0. Show directly from
your definition that if X = Z(f ), where f ∈ k[x1 , . . . , xn ] is irreducible, then dim X = n−1.
[Any form of the Nullstellensatz may be used if you state it clearly.]
(b) Suppose that char k = 0, and let W be the vector space of homogeneous
polynomials of degree d in 3 variables over k. Show that

U = {(f, p) ∈ W × k3 | Z(f − 1) is a smooth surface at p}

is a non-empty Zariski open subset of W × k3 .

Paper 1, Section II
25I Algebraic Geometry
(a) Let k be an uncountable field, M ⊆ k[x1 , . . . , xn ] a maximal ideal and
A = k[x1 , . . . , xn ]/M.
Show that every element of A is algebraic over k.
(b) Now assume that k is algebraically closed. Suppose that J ⊂ k[x1 , . . . , xn ] is
an ideal, and that f ∈ k[x1 , . . . , xn ] vanishes on Z(J). Using the result of part (a) or
otherwise, show that f N ∈ J for some N > 1.
(c) Let f : X → Y be a morphism of affine algebraic varieties. Show f (X) = Y if
and only if the map f ∗ : k[Y ] → k[X] is injective.
Suppose now that f (X) = Y , and that X and Y are irreducible. Define the
dimension of X, dim X, and show dim X > dim Y . [You may use whichever definition of
dim X you find most convenient.]

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Paper 3, Section II
20H Algebraic Topology
(a) State a version of the Seifert–van Kampen theorem for a cell complex X written
as the union of two subcomplexes Y, Z.
(b) Let
. . ∨ S }1 ∨ RP 2
Xn = |S 1 ∨ .{z
n

for n > 1, and take any x0 ∈ Xn . Write down a presentation for π1 (Xn , x0 ).
(c) By computing a homology group of a suitable four-sheeted covering space of
Xn , prove that Xn is not homotopy equivalent to a compact, connected surface whenever
n > 1.

Paper 2, Section II
21H Algebraic Topology
(a) Define the first barycentric subdivision K ′ of a simplicial complex K. Hence
define the r th barycentric subdivision K (r) . [You do not need to prove that K ′ is a simplicial
complex.]
(b) Define the mesh µ(K) of a simplicial complex K. State a result that describes
the behaviour of µ(K (r) ) as r → ∞.
(c) Define a simplicial approximation to a continuous map of polyhedra

f : |K| → |L|.

Prove that, if g is a simplicial approximation to f , then the realisation |g| : |K| → |L| is
homotopic to f .
(d) State and prove the simplicial approximation theorem. [You may use the
Lebesgue number lemma without proof, as long as you state it clearly.]
(e) Prove that every continuous map of spheres S n → S m is homotopic to a constant
map when n < m.

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Paper 1, Section II
21H Algebraic Topology
(a) Let V be the vector space of 3-dimensional upper-triangular matrices with real
entries:   
 1 x y 
V =  0 1 z  x, y, z ∈ R .
 
0 0 1
Let Γ be the set of elements of V for which x, y, z are integers. Notice that Γ is a subgroup
of GL3 (R); let Γ act on V by left-multiplication and let N = Γ\V . Show that the quotient
map V → N is a covering map.
(b) Consider the unit circle S 1 ⊆ C, and let T = S 1 × S 1 . Show that the map
f : T → T defined by
f (z, w) = (zw, w)
is a homeomorphism.
(c) Let M = [0, 1] × T / ∼, where ∼ is the smallest equivalence relation satisfying

(1, x) ∼ (0, f (x))

for all x ∈ T . Prove that N and M are homeomorphic by exhibiting a homeomor-

phism M → N . [You may assume without proof that N is Hausdorff.]
(d) Prove that π1 (M ) ∼
= Γ.

Paper 4, Section II
21H Algebraic Topology
(a) State the Mayer–Vietoris theorem for a union of simplicial complexes

K =M ∪N

with L = M ∩ N .
(b) Construct the map ∂∗ : Hk (K) → Hk−1 (L) that appears in the statement of the
theorem. [You do not need to prove that the map is well defined, or a homomorphism.]
(c) Let K be a simplicial complex with |K| homeomorphic to the n-dimensional
sphere S n , for n > 2. Let M ⊆ K be a subcomplex with |M | homeomorphic to
S n−1 × [−1, 1]. Suppose that K = M ∪ N , such that L = M ∩ N has polyhedron |L|
identified with S n−1 × {−1, 1} ⊆ S n−1 × [−1, 1]. Prove that |N | has two path components.

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Paper 3, Section II
22F Analysis of Functions
(a) Let (X, A, µ) be a measure space. Define the spaces Lp (X) for p ∈ [1, ∞]. Prove
that if µ(X) < ∞ then Lq (X) ⊂ Lp (X) for all 1 6 p < q 6 ∞.
(b) Now let X = Rn endowed with Borel sets and Lebesgue measure. Describe the
dual spaces of Lp (X) for p ∈ [1, ∞). Define reflexivity and say which Lp (X) are reflexive.
Prove that L1 (X) is not the dual space of L∞ (X).
(c) Now let X ⊂ Rn be a Borel subset and consider the measure space (X, A, µ)
induced from Borel sets and Lebesgue measure on Rn .

(i) Given any p ∈ [1, ∞], prove that any sequence (fn ) in Lp (X) converging
in Lp (X) to some f ∈ Lp (X) admits a subsequence converging almost
everywhere to f .

(ii) Prove that if Lq (X) ⊂ Lp (X) for 1 6 p < q 6 ∞ then µ(X) < ∞. [Hint:
You might want to prove first that the inclusion is continuous with the help
of one of the corollaries of Baire’s category theorem.]

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Paper 4, Section II
23F Analysis of Functions R
Here and below, Φ : R → R is smooth such that R e−Φ(x) dx = 1 and
!
|Φ′ (x)|2 Φ′′ (x)
lim − = ℓ ∈ (0, +∞).
|x|→+∞ 4 2

Cc1 (R) denotes the set of continuously differentiable complex-valued functions with com-
pact support on R.
(a) Prove that there are constants R0 > 0, λ1 > 0 and K1 > 0 so that for any
R > R0 and h ∈ Cc1 (R):
Z Z Z
2 −Φ(x) 2 −Φ(x)
′
h (x) e dx > λ1 |h(x)| e dx − K1 |h(x)|2 e−Φ(x) dx.
R {|x|>R} {|x|6R}

[Hint: Denote g := he−Φ/2 , expand the square and integrate by parts.]

(b) Prove that, given any R > 0, there is a CR > 0 so that for any h ∈ C 1 ([−R, R])
R +R
with −R h(x)e−Φ(x) dx = 0:
Z +R 1/2
|h(x) − h(y)| ′ 2 −Φ(x)
max |h(x)| + sup 1/2
6 CR |h (x)| e dx .
x∈[−R,R] {x,y∈[−R,R], x6=y} |x − y| −R

[Hint: Use the fundamental theorem of calculus to control the second term of the left-hand
side, and then compare h to its weighted mean to control the first term of the left-hand
side.]
(c) Prove that, given any R > 0, there is a λR > 0 so that for any h ∈ C 1 ([−R, R]):

Z Z R +R −Φ(y) dy 2
+R +R
′ 2 −Φ(x) −R h(y)e
h (x) e dx > λR h(x) − R +R e−Φ(x) dx.
−R −R e−Φ(y) dy
−R
R +R
[Hint: Show first that one can reduce to the case −R he−Φ = 0. Then argue by
contradiction with the help of the Arzelà–Ascoli theorem and part (b).]
(d) Deduce that there is a λ0 > 0 so that for any h ∈ Cc1 (R):
Z Z Z  2
′ 2 −Φ(x) −Φ(y)
h (x) e dx > λ0 h(x) − h(y)e dy e−Φ(x) dx.
R R R
R
[Hint: Show first that one can reduce to the case R he−Φ = 0. Then combine the inequality
(a), multiplied by a constant of the form ǫ = ǫ0 λR (where ǫ0 > 0 is chosen so that ǫ be suffi-
ciently small), and the inequality (c).]

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Paper 1, Section II
23F Analysis of Functions
(a) Consider a measure space (X, A, µ) and a complex-valued measurable function
F on X. Prove that for any ϕ : [0, +∞) → [0, +∞) differentiable and increasing such that
ϕ(0) = 0, then
Z Z +∞
ϕ(|F (x)|) dµ(x) = ϕ′ (s)µ({|F | > s}) dλ(s)
X 0

where λ is the Lebesgue measure.

(b) Consider a complex-valued measurable
R function f ∈ L1 (Rn ) ∩ L∞ (Rn ) and its
1
maximal function M f (x) = supr>0 |B(x,r)| B(x,r) |f | dλ. Prove that for p ∈ (1, +∞) there
is a constant cp > 0 such that kM f kLp (Rn ) 6 cp kf kLp (Rn ) .
[Hint: Split f = f0 + f1 with f0 = f χ{|f |>s/2} and f1 = f χ{|f |6s/2} and prove that
λ({M f > s}) 6 λ({M f0 > s/2}). Then use the maximal inequality λ({M f > s}) 6
C1
s kf kL1 (Rn ) for some constant C1 > 0.]
(c) Consider p, q ∈ (1, +∞) with p < q and α ∈ (0, n) such that 1/q = 1/p − α/n.
R |f (y)| αp/n 1−αp/n .
Define Iα |f |(x) := Rn |x−y| n−α dλ(y) and prove Iα |f |(x) 6 kf kLp (Rn ) M f (x)

[Hint: Split the integral into |x − y| > r and |x − y| ∈ [2−k−1 r, 2−k r) for all k > 0,
given some suitable r > 0.]

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Paper 1, Section II
34A Applications of Quantum Mechanics
A particle of mass m moves in one dimension in a periodic potential V (x) satisfying
V (x + a) = V (x). Define the Floquet matrix F . Show that det F = 1 and explain why
Tr F is real. Show that allowed bands occur for energies such that (Tr F )2 < 4. Consider
the potential
+∞
~2 λ X
V (x) = − δ(x − na).
m n=−∞

For states of negative energy, construct the Floquet matrix with respect to the basis of
states e±µx . Derive an inequality for the values of µ in an allowed energy band.
For states of positive energy, construct the Floquet matrix with respect to the basis
of states e±ikx . Derive an inequality for the values of k in an allowed energy band.
Show that the state with zero energy lies in a forbidden region for λa > 2.

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Paper 4, Section II
34A Applications of Quantum Mechanics
Define a Bravais lattice Λ in three dimensions. Define the reciprocal lattice Λ⋆ .
Define the Brillouin zone.
An FCC lattice has a basis of primitive vectors given by
a a a
a1 = (e2 + e3 ) , a2 = (e1 + e3 ) , a3 = (e1 + e2 ),
2 2 2
where ei is an orthonormal basis of R3 . Find a basis of reciprocal lattice vectors. What is
the volume of the Brillouin zone?
The asymptotic wavefunction for a particle, of wavevector k, scattering off a
potential V (r) is
eikr
ψ(r) ∼ eik·r + fV (k; k′ ) ,
r
where k′ = kr̂ and fV (k; k′ ) is the scattering amplitude. Give a formula for the Born
approximation to the scattering amplitude.
Scattering of a particle off a single atom is modelled by a potential V (r) = V0 δ(r−d)
with δ-function support on a spherical shell, r = |r| = d centred at the origin. Calculate
the Born approximation to the scattering amplitude, denoting the resulting expression as
f˜V (k; k′ ).
Scattering of a particle off a crystal consisting of atoms located at the vertices of a
lattice Λ is modelled by a potential
X
VΛ = V (r − R),
R∈Λ

where V (r) = V0 δ(r − d) as above. Calculate the Born approximation to the scattering
amplitude giving your answer in terms of your approximate expression f˜V for scattering off
a single atom. Show that the resulting amplitude vanishes unless the momentum transfer
q = k − k′ lies in the reciprocal lattice Λ⋆ .
For the particular FCC lattice considered above, show that, when k = |k| > 2π/a,
scattering occurs for two values of the scattering angle, θ1 and θ2 , related by
 
sin θ21 2
 =√ .
sin θ2 3
2

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Paper 3, Section II
35A Applications of Quantum Mechanics
A beam of particles of mass m and momentum p = ~k is incident along the z-axis.
The beam scatters off a spherically symmetric potential V (r). Write down the asymptotic
form of the wavefunction in terms of the scattering amplitude f (θ).
The incoming plane wave and the scattering amplitude can be expanded in partial
waves as,
1 X
∞  
eikr cos θ ∼ (2l + 1) eikr − (−1)l e−ikr Pl (cos θ)
2ikr
l=0
∞
X 2l + 1
f (θ) = fl Pl (cos θ)
k
l=0

where Pl are Legendre polynomials. Define the S-matrix. Assuming that the S-matrix is
unitary, explain why we can write

fl = eiδl sin δl

for some real phase shifts δl . Obtain an expression for the total cross-section σT in terms
of the phase shifts δl .
[Hint: You may use the orthogonality of Legendre polynomials:
Z +1
2
dw Pl (w)Pl′ (w) = δll′ . ]
−1 2l + 1

Consider the repulsive, spherical potential


+V0 r<a
V (r) =
0 r>a

where V0 = ~2 γ 2 /2m. By considering the s-wave solution to the Schrödinger equation,

show that p
tan(ka + δ0 ) tanh( γ 2 − k2 a)
= p .
ka γ 2 − k2 a
For low momenta, ka ≪ 1, compute the s-wave contribution to the total cross-section.
Comment on the physical interpretation of your result in the limit γa → ∞.

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Paper 2, Section II
35A Applications of Quantum Mechanics
Consider a one-dimensional chain of 2N ≫ 1 atoms, each of mass m. Impose
periodic boundary conditions. The forces between neighbouring atoms are modelled
as springs, with alternating spring constants λ and αλ. In equilibrium, the separation
between the atoms is a.
Denote the position of the nth atom as xn (t). Let un (t) = xn (t) − na be the
displacement from equilibrium. Write down the equations of motion of the system.
Show that the longitudinal modes of vibration are labelled by a wavenumber k
that is restricted to lie in a Brillouin zone. Find the frequency spectrum. What is the
frequency gap at the edge of the Brillouin zone? Show that the gap vanishes when α = 1.
Determine approximations for the frequencies near the centre of the Brillouin zone. Plot
the frequency spectrum. What is the speed of sound in this system?

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Paper 4, Section II
27J Applied Probability
Let X1 , X2 , . . . be independent, identically distributed random variables with finite
mean µ. Explain what is meant by saying that the random variable M is a stopping time
with respect to the sequence (Xi : i = 1, 2, . . . ).
Let M be a stopping time with finite mean E(M ). Prove Wald’s equation:
M
!
X
E Xi = µE(M ).
i=1

[Here and in the following, you may use any standard theorem about integration.]
Suppose the Xi are strictly positive, and let N be the renewal process with
interarrival times (Xi : i = 1, 2, . . . ). Prove that m(t) = E(Nt ) satisfies the elementary
renewal theorem:
1 1
m(t) → as t → ∞.
t µ

A computer keyboard contains 100 different keys, including the lower and upper
case letters, the usual symbols, and the space bar. A monkey taps the keys uniformly at
random. Find the mean number of keys tapped until the first appearance of the sequence
‘lava’ as a sequence of 4 consecutive characters.
Find the mean number of keys tapped until the first appearance of the sequence
‘aa’ as a sequence of 2 consecutive characters.

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Paper 3, Section II
27J Applied Probability
Individuals arrive in a shop in the manner of a Poisson process with intensity λ,
and they await service in the order of their arrival. Their service times are independent,
identically distributed random variables S1 , S2 , . . . . For n > 1, let Qn be the number
remaining in the shop immediately after the nth departure. Show that

Qn+1 = An + Qn − h(Qn ),

where An is the number of arrivals during the (n + 1)th service period, and h(x) =
min{1, x}.
Show that
E(An ) = ρ, E(A2n ) = ρ + λ2 E(S 2 ),
where S is a typical service period, and ρ = λE(S) is the traffic intensity of the queue.
Suppose ρ < 1, and the queue is in equilibrium in the sense that Qn and Qn+1 have
the same distribution for all n. Express E(Qn ) in terms of λ, E(S), E(S 2 ). Deduce that
the mean waiting time (prior to service) of a typical individual is 12 λE(S 2 )/(1 − ρ).

Paper 2, Section II
27J Applied Probability
Let X = (Xt : t > 0) be a continuous-time Markov chain on the finite state space
S. Define the terms generator (or Q-matrix ) and invariant distribution, and derive an
equation that links the generator G and any invariant distribution π. Comment on the
possible non-uniqueness of invariant distributions.
Suppose X is irreducible, and let N = (Nt : t > 0) be a Poisson process with
intensity λ, that is independent of X. Let Yn be the value of X immediately after the nth
arrival-time of N (and Y0 = X0 ). Show that (Yn : n = 0, 1, . . .) is a discrete-time Markov
chain, state its transition matrix and prove that it has the same invariant distribution
as X.

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Paper 1, Section II
28J Applied Probability
Let λ : [0, ∞) → (0, ∞) be a continuous function. Explain what is meant by an
inhomogeneous Poisson process with intensity function λ.
Let (Nt : t > 0) be such an inhomogeneous Poisson process, and let Mt = Ng(t)
where g : [0, ∞) → [0, ∞) is strictly increasing, differentiable and satisfies g(0) = 0. Show
that MRis a homogeneous Poisson process with intensity 1 if Λ(g(t)) = t for all t, where
t
Λ(t) = 0 λ(u) du. Deduce that Nt has the Poisson distribution with mean Λ(t).
Bicycles arrive at the start of a long road in the manner of a Poisson process
N = (Nt : t > 0) with constant intensity λ. The ith bicycle has constant velocity
Vi , where V1 , V2 , . . . are independent, identically distributed random variables, which are
independent of N . Cyclists can overtake one another freely. Show that the number
of bicycles on the first x miles
of the road at time t has the Poisson distribution with
parameter λE V −1 min{x, V t} .

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Paper 2, Section II
31B Asymptotic Methods
R +∞ R +R
2 √ 2
Given that −∞ e−u du = π obtain the value of limR→+∞ −R e−itu du for real
RR 3
positive t. Also obtain the value of limR→+∞ 0 e−itu du, for real positive t, in terms of
R +∞ 3
Γ( 34 ) = 0 e−u du.
For α > 0, x > 0, let
Z π
1

Qα (x) = cos x sin θ − αθ dθ .
π 0

Find the leading terms in the asymptotic expansions as x → +∞ of (i) Qα (x) with α fixed,
and (ii) of Qx (x).

Paper 3, Section II
31B Asymptotic Methods
(a) Find the curves of steepest descent emanating from t = 0 for the integral
Z
1
Jx (x) = ex(sinh t−t) dt ,
2πi C

for x > 0 and determine the angles at which they meet at t = 0, and their asymptotes at
infinity.
(b) An integral representation for the Bessel function Kν (x) for real x > 0 is
Z  
1 +∞ νh(t) x
Kν (x) = e dt , h(t) = t − cosh t .
2 −∞ ν

Show that, as ν → +∞ , with x fixed,

 1  ν
π 2 2ν
Kν (x) ∼ .
2ν ex

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Paper 4, Section II
31B Asymptotic Methods
Show that Z π
1
I0 (x) = ex cos θ dθ
π 0
is a solution to the equation
xy ′′ + y ′ − xy = 0 ,
and obtain the first two terms in the asymptotic expansion of I0 (x) as x → +∞.
1
For x > 0, define a new dependent variable w(x) = x 2 y(x), and show that if y
solves the preceding equation then
 
′′ 1
w + − 1 w = 0.
4x2

Obtain the Liouville–Green approximate solutions to this equation for large positive x,
and compare with your asymptotic expansion for I0 (x) at the leading order.

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Paper 4, Section I
4G Automata and Formal Languages
(a) State the s-m-n theorem, the recursion theorem, and Rice’s theorem.
(b) Show that if g : N2 → N is partial recursive, then there is some e ∈ N such that

fe,1 (y) = g(e, y) ∀y ∈ N.

(c) By considering the partial function g : N2 → N given by


0 if y < x
g(x, y) =
↑ otherwise,

show there exists some m ∈ N such that Wm has exactly m elements.

(d) Given n ∈ N, is it possible to compute whether or not Wn has exactly 9 elements?
Justify your answer.
[Note that we define N = {0, 1, . . .}. Any use of Church’s thesis in your answers
should be explicitly stated.]

Paper 3, Section I
4G Automata and Formal Languages
(a) Define what it means for a context-free grammar (CFG) to be in Chomsky
normal form (CNF).
(b) Give an algorithm for converting a CFG G into a corresponding CFG GChom in
CNF satisfying L(GChom ) = L(G) − {ǫ}. [You need only outline the steps, without proof.]
(c) Convert the following CFG G:

S → ASc | B , A → a , B → b ,

into a grammar in CNF.

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Paper 2, Section I
4G Automata and Formal Languages
(a) Let E = (QE , Σ, δE , q0 , FE ) be a nondeterministic finite-state automaton with
ǫ-transitions (ǫ-NFA). Define the deterministic finite-state automaton (DFA)
D = (QD , Σ, δD , qD , FD ) obtained from E via the subset construction with
ǫ-transitions.
(b) Let E and D be as above. By inducting on lengths of words, prove that

δ̂E (q0 , w) = δ̂D (qD , w) for all w ∈ Σ∗ .

(c) Deduce that L(D) = L(E).

Paper 1, Section I
4G Automata and Formal Languages
(a) State the pumping lemma for context-free languages (CFLs).
(b) Which of the following are CFLs? Justify your answers.

(i) {ww | w ∈ {a, b, c}∗ }

(ii) {am bn ck dl | 3m = 4n and 2k = 5l}

n
(iii) {a3 | n > 0}

(c) Let L be a CFL. Show that L∗ is also a CFL.

Paper 3, Section II
12G Automata and Formal Languages
(a) State and prove the pumping lemma for regular languages.
(b) Let D be a minimal deterministic finite-state automaton whose language L(D)
is finite. Let ΓD be the transition diagram of D, and suppose there exists a non-empty
closed path γ in ΓD starting and ending at state p.

(i) Show that there is no path in ΓD from p to any accept state of D.

(ii) Show that there is no path in ΓD from p to any other state of D.

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Paper 1, Section II
12G Automata and Formal Languages
(a) Define the halting set K. Prove that K is recursively enumerable, but not
recursive.
(b) Given A, B ⊆ N, define a many-one reduction of A to B. Show that if B is
recursively enumerable and A 6m B, then A is also recursively enumerable.
(c) Show that each of the functions f (n) = 2n and g(n) = 2n + 1 are both partial
recursive and total, by building them up as partial recursive functions.
(d) Let X, Y ⊆ N. We define the set X ⊕ Y via

X ⊕ Y := {2x | x ∈ X} ∪ {2y + 1 | y ∈ Y }.

(i) Show that both X 6m X ⊕ Y and Y 6m X ⊕ Y .

(ii) Using the above, or otherwise, give an explicit example of a subset C of N

for which neither C nor N \ C are recursively enumerable.

(iii) For every Z ⊆ N, show that if X 6m Z and Y 6m Z then X ⊕ Y 6m Z.

[Note that we define N = {0, 1, . . .}. Any use of Church’s thesis in your answers
should be explicitly stated.]

Part II, 2018 List of Questions [TURN OVER

 22

Paper 1, Section I
8B Classical Dynamics
Derive Hamilton’s equations from an action principle.
Consider a two-dimensional phase space with the Hamiltonian H = p2 + q −2 . Show
that F = pq − ctH is the first integral for some constant c which should be determined.
By considering the surfaces of constant F in the extended phase space, solve Hamilton’s
equations, and sketch the orbits in the phase space.

Paper 2, Section I
8B Classical Dynamics
Let x = xi + yj + zk. Consider a Lagrangian
1 2
L= ẋ + y ẋ
2
of a particle constrained to move on a sphere |x| = 1/c of radius 1/c. Use Lagrange
multipliers to show that

ẍ + ẏi − ẋj + c2 (|ẋ|2 + y ẋ − xẏ)x = 0. (∗)

Now, consider the system (∗) with c = 0, and find the particle trajectories.

Paper 3, Section I
8B Classical Dynamics
Three particles of unit mass move along a line in a potential
1 
V = (x1 − x2 )2 + (x1 − x3 )2 + (x3 − x2 )2 + x21 + x22 + x23 ,
2
where xi is the coordinate of the i’th particle, i = 1, 2, 3.
Write the Lagrangian in the form
1 1
L = Tij ẋi ẋj − Vij xi xj ,
2 2
and specify the matrices Tij and Vij .
Find the normal frequencies and normal modes for this system.

Part II, 2018 List of Questions

 23

Paper 4, Section I
8B Classical Dynamics
State and prove Noether’s theorem in Lagrangian mechanics.
Consider a Lagrangian
 
1 ẋ2 + ẏ 2 x
L= −V
2 y2 y

for a particle moving in the upper half-plane {(x, y) ∈ R2 , y > 0} in a potential V which
only depends on x/y. Find two independent first integrals.

Paper 2, Section II
14B Classical Dynamics
Define a body frame ea (t), a = 1, 2, 3 of a rotating rigid body, and show that there
exists a vector ω = (ω1 , ω2 , ω3 ) such that

ėa = ω × ea .

Let L = I1 ω1 (t)e1 + I2 ω2 (t)e2 + I3 ω3 (t)e3 be the angular momentum of a free rigid

body expressed in the body frame. Derive the Euler equations from the conservation of
angular momentum.
Verify that the kinetic energy E, and the total angular momentum L2 are conserved.
Hence show that
ω̇32 = f (ω3 ),
where f (ω3 ) is a quartic polynomial which should be explicitly determined in terms of L2
and E.

Paper 4, Section II
15B Classical Dynamics
Given a Lagrangian L(qi , q̇i , t) with degrees of freedom qi , define the Hamiltonian
and show how Hamilton’s equations arise from the Lagrange equations and the Legendre
transform.
Consider the Lagrangian for a symmetric top moving in constant gravity:
1 1
L= A(θ̇ 2 + φ̇2 sin2 θ) + B(ψ̇ + φ̇ cos θ)2 − M gl cos θ,
2 2
where A, B, M , g and l are constants. Construct the corresponding Hamiltonian, and
find three independent Poisson-commuting first integrals of Hamilton’s equations.

Part II, 2018 List of Questions [TURN OVER

 24

Paper 4, Section I
3H Coding & Cryptography
What is a linear feedback shift register ? Explain the Berlekamp–Massey method
for recovering a feedback polynomial of a linear feedback shift register from its output.
Illustrate the method in the case when we observe output

0 1 0 1 1 1 1 0 0 0 1 0....

Paper 3, Section I
3H Coding & Cryptography
Compute the rank and minimum distance of the cyclic code with generator polyno-
mial g(X) = X 3 + X 2 + 1 and parity check polynomial h(X) = X 4 + X 3 + X 2 + 1. Now let
α be a root of g(X) in the field with 8 elements. We receive the word r(X) = X 2 + X + 1
(mod X 7 − 1). Verify that r(α) = α4 , and hence decode r(X) using minimum-distance
decoding.

Paper 2, Section I
3H Coding & Cryptography
What is the channel matrix of a binary symmetric channel with error probability p?
State the maximum likelihood decoding rule and the minimum distance decoding
rule. Prove that if p < 1/2, then they agree.
Let C be the repetition code {000, 111}. Suppose a codeword from C is sent through
a binary symmetric channel with error probability p. Show that, if the minimum distance
decoding rule is used, then the probability of error is 3p2 − 2p3 .

Paper 1, Section I
3H Coding & Cryptography
State and prove Shannon’s noiseless coding theorem. [You may use Gibbs’ and
Kraft’s inequalities as long as they are clearly stated.]

Part II, 2018 List of Questions

 25

Paper 1, Section II
11H Coding & Cryptography
Define the bar product C1 |C2 of binary linear codes C1 and C2 , where C2 is a subcode
of C1 . Relate the rank and minimum distance of C1 |C2 to those of C1 and C2 and justify
your answer.
What is a parity check matrix for a linear code? If C1 has parity check matrix P1
and C2 has parity check matrix P2 , find a parity check matrix for C1 |C2 .
Using the bar product construction, or otherwise, define the Reed–Muller code
RM (d, r) for 0 6 r 6 d. Compute the rank of RM (d, r). Show that all but two codewords
in RM (d, 1) have the same weight. Given d, for which r is it true that all elements of
RM (d, r) have even weight? Justify your answer.

Paper 2, Section II
12H Coding & Cryptography
Describe the RSA encryption scheme with public key (N, e) and private key d.
Suppose N = pq with p and q distinct odd primes and 1 6 x 6 N with x and N
coprime. Denote the order of x in F∗p by Op (x). Further suppose Φ(N ) divides 2a b where
b is odd. If Op (xb ) 6= Oq (xb ) prove that there exists 0 6 t < a such that the greatest
t
common divisor of x2 b − 1 and N is a nontrivial factor of N . Further, prove that the
number of x satisfying Op (xb ) 6= Oq (xb ) is > Φ(N )/2.
Hence, or otherwise, prove that finding the private key d from the public key (N, e)
is essentially as difficult as factoring N .
Suppose a message m is sent using the RSA scheme with e = 43 and N = 77, and
c = 5 is the received text. What is m?
An integer m satisfying 1 6 m 6 N − 1 is called a fixed point if it is encrypted to
itself. Prove that if m is a fixed point then so is N − m.

Part II, 2018 List of Questions [TURN OVER

 26

Paper 2, Section I
9B Cosmology
(a) Consider a homogeneous and isotropic universe with a uniform distribution of
galaxies. For three galaxies at positions rA , rB , rC , show that spatial homogeneity implies
that their non-relativistic velocities v(r) must satisfy

v(rB − rA ) = v(rB − rC ) − v(rA − rC ),

and hence that the velocity field coordinates vi are linearly related to the position
coordinates rj via

vi = Hij rj ,

where the matrix coefficients Hij are independent of the position. Show why isotropy then
implies Hubble’s law

v = H r, with H independent of r .

Explain how the velocity of a galaxy is determined by the scale factor a and express the
Hubble parameter H0 today in terms of a.
(b) Define the cosmological horizon dH (t). For an Einstein–de Sitter universe with
a(t) ∝ t2/3 , calculate dH (t0 ) at t = t0 today in terms of H0 . Briefly describe the horizon
problem of the standard cosmology.

Part II, 2018 List of Questions

 27

Paper 3, Section I
9B Cosmology
The energy density of a particle species is defined by
Z ∞
ǫ= E(p)n(p)dp ,
0
p
where E(p) = c p2 + m2 c2 is the energy, and n(p) the distribution function, of a particle
with momentum p. Here c is the speed of light and m is the rest mass of the particle. If
the particle species is in thermal equilibrium then the distribution function takes the form

4π p2
n(p) = g ,
h3 exp((E(p) − µ)/kT ) ∓ 1

where g is the number of degrees of freedom of the particle, T is the temperature, h and
k are constants and − is for bosons and + is for fermions.
(a) Stating any assumptions you require, show that in the very early universe the
energy density of a given particle species i is
Z ∞
4πgi 4 y3
ǫi = (kT ) dy .
(hc)3 0 ey ∓ 1

(b) Show that the total energy density in the very early universe is

4π 5
ǫ= g∗ (kT )4 ,
15(hc)3

where g∗ is defined by
X 7 X
g∗ ≡ gi + gi .
8
Bosons Fermions
R∞ 3 y
[Hint: You may use the fact that 0 y (e − 1)−1 dy = π 4 /15.]

Part II, 2018 List of Questions [TURN OVER

 28

Paper 1, Section I
9B Cosmology
For a homogeneous and isotropic universe filled with pressure-free matter (P = 0),
the Friedmann and Raychaudhuri equations are, respectively,
 2
ȧ kc2 8πG ä 4πG
+ 2 = ρ and =− ρ,
a a 3 a 3

with mass density ρ, curvature k, and where ȧ ≡ da/dt. Using conformal time τ with
dτ = dt/a, show that the relative density parameter can be expressed as

ρ 8πGρ a2
Ω(t) ≡ = ,
ρcrit 3H2
da
where H = a1 dτ and ρcrit is the critical density of a flat k = 0 universe (Einstein–de Sitter).
Use conformal time τ again to show that the Friedmann and Raychaudhuri equations can
be re-expressed as

kc2 dH
= Ω−1 and 2 + H2 + kc2 = 0 .
H2 dτ
From these derive the evolution equation for the density parameter Ω:
dΩ
= H Ω (Ω − 1) .
dτ
Plot the qualitative behaviour of Ω as a function of time relative to the expanding Einstein–
de Sitter model with Ω = 1 (i.e., include curves initially with Ω > 1 and Ω < 1).

Part II, 2018 List of Questions

 29

Paper 4, Section I
9B Cosmology
A constant overdensity is created by taking a spherical region of a flat matter-
dominated universe with radius R and compressing it into a region with radius R < R.
The evolution is governed by the parametric equations

R = AR0 (1 − cos θ) , t = B(θ − sin θ),

where R0 is a constant and

Ωm,0 Ωm,0
A= , B= ,
2(Ωm,0 − 1) 2H0 (Ωm,0 − 1)3/2

where H0 is the Hubble constant and Ωm,0 is the fractional overdensity at time t0 .
Show that, as t → 0+ ,
 2/3 !
1/3 1 6t
R(t) = R0 Ωm,0 a(t) 1− + ... ,
20 B

where the scale factor is given by a(t) = (3H0 t/2)2/3 .

Show that, at the linear level, the density perturbation δlinear grows as a(t). Show
that, when the spherical overdensity has collapsed to zero radius, the linear perturbation
3
has value δlinear = 20 (12π)2/3 .

Part II, 2018 List of Questions [TURN OVER

 30

Paper 3, Section II
14B Cosmology
The pressure support equation for stars is
 
1 d r 2 dP
= −4πGρ ,
r 2 dr ρ dr

where ρ is the density, P is the pressure, r is the radial distance, and G is Newton’s
constant.
(a) What two boundary conditions should we impose on the above equation for it
to describe a star?
(b) By assuming a polytropic equation of state,
1
P (r) = Kρ1+ n (r) ,

where K is a constant, derive the Lane–Emden equation

 
1 d 2 dθ
ξ = −θ n ,
ξ 2 dξ dξ

where ρ = ρc θ n , with ρc the density at the centre of the star, and r = aξ, for some a that
you should determine.
(c) Show that the mass of a polytropic star is
 3
1 (n + 1) K 2 3−n
M= √ ρc 2n Yn ,
2 π G

where Yn ≡ − ξ12 dθ
dξ and ξ1 is the value of ξ at the surface of the star.
ξ=ξ1

(d) Derive the following relation between the mass, M , and radius, R, of a polytropic
star
n 3−n
M = An K n−1 R 1−n ,

where you should determine the constant An . What type of star does the n = 3 polytrope
represent and what is the significance of the mass being constant for this star?

Part II, 2018 List of Questions

 31

Paper 1, Section II
15B Cosmology
A flat (k=0) homogeneous and isotropic universe with scale factor a(t) is filled with
a scalar field φ(t) with potential V (φ). Its evolution satisfies the Friedmann and scalar
field equations,
 
1 1 2 dV
H2 = 2 φ̇ + c2
V (φ) , φ̈ + 3H φ̇ + c2 = 0,
3MPl 2 dφ

where H(t) = ȧa is the Hubble parameter, MPl is the reduced Planck mass, and dots denote
derivatives with respect to cosmic time t, e.g. φ̇ ≡ dφ/dt.
(a) Use these equations to derive the Raychaudhuri equation, expressed in the form:
1 2
Ḣ = − 2 φ̇ .
2MPl

(b) Consider the following ansatz for the scalar field evolution,

φ(t) = φ0 ln tanh(λt) , (†)

where λ, φ0 are constants. Find the specific cosmological solution,

φ20
H(t) = λ 2 coth(2λt) ,
MPl
2 2
a(t) = a0 [sinh(2λt)]φ0 /2MPl , a0 constant.

(c) Hence, show that the Hubble parameter can be expressed in terms of φ as
 
φ20 φ
H(φ) = λ 2 cosh ,
MPl φ0

and that the scalar field ansatz solution (†) requires the following form for the potential:
    
2λ2 φ20 3φ20 2 φ
V (φ) =
c2 2MPl2 − 1 cosh φ0
+1 .

(d) Assume that the given parameters in V (φ) are such that 2/3 < φ20 /MPl 2 < 2.

Show that the asymptotic limit for the cosmological solution as t → 0 exhibits decelerating
power law evolution and that there is an accelerating solution as t → ∞, that is,
2 2
t → 0, φ → −∞ , a(t) ∼ tφ0 /2MPl ,
t → ∞, φ → 0, a(t) ∼ exp(λφ20 t/MPl
2
).

Find the time tacc at which the solution transitions from deceleration to acceleration.

Part II, 2018 List of Questions [TURN OVER

 32

Paper 4, Section II
25I Differential Geometry
Let S ⊂ R3 be a surface.
(a) Define what it means for a curve γ : I → S to be a geodesic, where I = (a, b)
and −∞ 6 a < b 6 ∞.
e : Ie → S with I ⊂ Ie
(b) A geodesic γ : I → S is said to be maximal if any geodesic γ
and γ e
e|I = γ satisfies I = I. A surface is said to be geodesically complete if all maximal
geodesics are defined on I = (−∞, ∞), otherwise, the surface is said to be geodesically
incomplete. Give an example, with justification, of a non-compact geodesically complete
surface S which is not a plane.
(c) Assume that along any maximal geodesic

γ : (−T− , T+ ) → S,

the following holds:

T± < ∞ =⇒ lim sup |K(γ(±s))| = ∞. (∗)
s→T±

Here K denotes the Gaussian curvature of S.

(i) Show that S is inextendible, i.e. if Se ⊂ R3 is a connected surface with

e then Se = S.
S ⊂ S,

(ii) Give an example of a surface S which is geodesically incomplete and sat-

isfies (∗). Do all geodesically incomplete inextendible surfaces satisfy (∗)?
Justify your answer.

[You may use facts about geodesics from the course provided they are clearly stated.]

Part II, 2018 List of Questions

 33

Paper 3, Section II
25I Differential Geometry
Let S ⊂ R3 be a surface.
(a) Define the Gaussian curvature K of S in terms of the coefficients of the first
and second fundamental forms, computed with respect to a local parametrization φ(u, v)
of S.
Prove the Theorema Egregium, i.e. show that the Gaussian curvature can be
expressed entirely in terms of the coefficients of the first fundamental form and their
first and second derivatives with respect to u and v.
(b) State the global Gauss–Bonnet theorem for a compact orientable surface S.
(c) Now assume that S is non-compact and diffeomorphic to S2 \ {(1, 0, 0)} but that
there is a Rpoint p ∈ R3 such that S ∪ {p} is a compact subset of R3 . Is it necessarily the
case that S KdA = 4/π? Justify your answer.

Paper 2, Section II
25I Differential Geometry
Let γ(t) : [a, b] → R3 denote a regular curve.
(a) Show that there exists a parametrization of γ by arc length.
(b) Under the assumption that the curvature is non-zero, define the torsion of γ.
Give an example of two curves γ1 and γ2 in R3 whose curvature (as a function of arc
length s) coincides and is non-vanishing, but for which the curves are not related by a
rigid motion, i.e. such that γ1 (s) is not identically ρ(R,T ) (γ2 (s)) where R ∈ SO(3), T ∈ R3
and
ρ(R,T ) (v) := T + Rv.

(c) Give an example of a simple closed curve γ, other than a circle, which is preserved
by a non-trivial rigid motion, i.e. which satisfies

ρ(R,T ) (v) ∈ γ([a, b]) for all v ∈ γ([a, b])

for some choice of R ∈ SO(3), T ∈ R3 with (R, T ) 6= (Id, 0). Justify your answer.
(d) Now show that a simple closed curve γ which is preserved by a nontrivial smooth
1-parameter family of rigid motions is necessarily a circle, i.e. show the following:
Let (R, T ) : (−ǫ, ǫ) → SO(3) × R3 be a regular curve. If for all t̃ ∈ (−ǫ, ǫ),

ρ(R(t̃),T (t̃)) (v) ∈ γ([a, b]) for all v ∈ γ([a, b]),

then γ([a, b]) is a circle. [You may use the fact that the set of fixed points of a non-trivial
rigid motion is either ∅ or a line L ⊂ R3 .]

Part II, 2018 List of Questions [TURN OVER

 34

Paper 1, Section II
26I Differential Geometry
(a) Let X ⊂ Rn be a manifold and p ∈ X. Define the tangent space Tp X and show
that it is a vector subspace of Rn , independent of local parametrization, of dimension
equal to dim X.
(b) Now show that Tp X depends continuously on p in the following sense: if pi is a
sequence in X such that pi → p ∈ X, and wi ∈ Tpi X is a sequence such that wi → w ∈ Rn ,
then w ∈ Tp X. If dim X > 0, show that all w ∈ Tp X arise as such limits where pi is a
sequence in X \ p.
(c) Consider the set Xa ⊂ R4 defined by Xa = {x21 + 2x22 = a2 } ∩ {x3 = ax4 }, where
a ∈ R. Show that, for all a ∈ R, the set Xa is a smooth manifold. Compute its dimension.
(d) For Xa as above, does Tp Xa depend continuously on p and a for all a ∈ R? In
other words, let ai ∈ R, pi ∈ Xai be sequences with ai → a ∈ R, pi → p ∈ Xa . Suppose
that wi ∈ Tpi Xai and wi → w ∈ R4 . Is it necessarily the case that w ∈ Tp Xa ? Justify
your answer.

Part II, 2018 List of Questions

 35

Paper 1, Section II
31E Dynamical Systems
Consider the system

ẋ = −2ax + 2xy , ẏ = 1 − x2 − y 2 ,

where a is a constant.
(a) Find and classify the fixed points of the system. For a = 0 show that the linear
classification of the non-hyperbolic fixed points is nonlinearly correct. For a 6= 0 show that
there are no periodic orbits. [Standard results for periodic orbits may be quoted without
proof.]
1 3
(b) Sketch the phase plane for the cases (i) a = 0, (ii) a = 2, and (iii) a = 2,
showing any separatrices clearly.
(c) For what values of a do stationary bifurcations occur? Consider the bifurcation
with a > 0. Let y0 , a0 be the values of y, a at which the bifurcation occurs, and define
Y = y − y0 , µ = a − a0 . Assuming that µ = O(x2 ), find the extended centre manifold
Y = Y (x, µ) to leading order. Further, determine the evolution equation on the centre
manifold to leading order. Hence identify the type of bifurcation.

Paper 4, Section II
32E Dynamical Systems
Let F : I → I be a continuous one-dimensional map of an interval I ⊂ R. Define
what it means (i) for F to have a horseshoe (ii) for F to be chaotic. [Glendinning’s
definition should be used throughout this question.]
Prove that if F has a 3-cycle x1 < x2 < x3 then F is chaotic. [You may assume the
intermediate value theorem and any corollaries of it.]
State Sharkovsky’s theorem.
Use the above results to deduce that if F has an N -cycle, where N is any integer
that is not a power of 2, then F is chaotic.
Explain briefly why if F is chaotic then F has N -cycles for many values of N that
are not powers of 2. [You may assume that a map with a horseshoe acts on some set Λ
like the Bernoulli shift map acts on [0,1).]
The√ logistic map is not chaotic when µ < µ∞ ≈ 3.57 and it has 3-cycles when
µ > 1 + 8 ≈ 3.84. What can be deduced from these statements about the values of µ
for which the logistic map has a 10-cycle?

Part II, 2018 List of Questions [TURN OVER

 36

Paper 3, Section II
32E Dynamical Systems
Consider the system

ẋ = y , ẏ = µ1 x + µ2 y − (x + y)3 ,

where µ1 and µ2 are parameters.

By considering a function of the form V (x, y) = f (x + y) + 21 y 2 , show that when
µ1 = µ2 = 0 the origin is globally asymptotically stable. Sketch the phase plane for this
case.
Find the fixed points for the general case. Find the values of µ1 and µ2 for which
the fixed points have (i) a stationary bifurcation and (ii) oscillatory (Hopf) bifurcations.
Sketch these bifurcation values in the (µ1 , µ2 )-plane.
For the case µ2 = −1, find the leading-order approximation to the extended centre
manifold of the bifurcation as µ1 varies, assuming that µ1 = O(x2 ). Find also the evolution
equation on the extended centre manifold to leading order. Deduce the type of bifurcation,
and sketch the bifurcation diagram in the (µ1 , x)-plane.

Part II, 2018 List of Questions

 37

Paper 2, Section II
32E Dynamical Systems
Consider the system

ẋ = y, ẏ = x − x3 + ǫ(1 − αx2 )y ,

where α and ǫ are real constants, and 0 6 ǫ ≪ 1. Find and classify the fixed points.
Show that when ǫ = 0 the system is Hamiltonian and find H. Sketch the phase
plane for this case.
Suppose now that 0 < ǫ ≪ 1. Show that the small change in H following a trajectory
of the perturbed system around an orbit H = H0 of the unperturbed system is given to
leading order by an equation of the form
Z x2
∆H = ǫ F (x; α, H0 ) dx ,
x1

where F should be found explicitly, and where x1 and x2 are the minimum and maximum
values of x on the unperturbed orbit.
Use the energy-balance method to find the value of α, correct to leading order in
ǫ, for which the system has a homoclinic orbit. [Hint: The substitution u = 1 − 21 x2 may
prove useful.]
Over what range of α would you expect there to be periodic solutions that enclose
only one of the fixed points?

Part II, 2018 List of Questions [TURN OVER

 38

Paper 1, Section II
36D Electrodynamics
Define the field strength tensor F µν (x) for an electromagnetic field specified by
a 4-vector potential Aµ (x). How do the components of F µν change under a Lorentz
transformation? Write down two independent Lorentz-invariant quantities which are
quadratic in the field strength tensor.
[Hint: The alternating tensor εµνρσ takes the values +1 and −1 when (µ, ν, ρ, σ) is
an even or odd permutation of (0, 1, 2, 3) respectively and vanishes otherwise. You may
assume this is an invariant tensor of the Lorentz group. In other words, its components
are the same in all inertial frames.]
In an inertial frame with spacetime coordinates xµ = (ct, x), the 4-vector potential
has components Aµ = (φ/c, A) and the electric and magnetic fields are given as

∂A
E = −∇φ −
∂t
B = ∇ × A.

Evaluate the components of F µν in terms of the components of E and B. Show that the
quantities
1
S = |B|2 − |E|2 and T =E·B
c2
are the same in all inertial frames.
A relativistic particle of mass m, charge q and 4-velocity uµ (τ ) moves according to
the Lorentz force law,
duµ q
= F µν uν . (∗)
dτ m
Here τ is the proper time. For the case of a constant, uniform field, write down a solution
of (∗) giving uµ (τ ) in terms of its initial value uµ (0) as an infinite series in powers of the
field strength.
Suppose further that the fields are such that both S and T defined above are zero.
Work in an inertial frame with coordinates xµ = (ct, x, y, z) where the particle is at rest at
the origin at t = 0 and the magnetic field points in the positive z-direction with magnitude
|B| = B. The electric field obeys E · ŷ = 0. Show that the particle moves on the curve
y 2 = Ax3 for some constant A which you should determine.

Part II, 2018 List of Questions

 39

Paper 4, Section II
36D Electrodynamics
(a) Define the polarisation of a dielectric material and explain what is meant by the
term bound charge.
Consider a sample of material with spatially dependent polarisation P(x) occupying
a region V with surface S. Show that, in the absence of free charge, the resulting scalar
potential φ(x) can be ascribed to bulk and surface densities of bound charge.
Consider a sphere of radius R consisting of a dielectric material with permittivity
ǫ surrounded by a region of vacuum. A point-like electric charge q is placed at the centre
of the sphere. Determine the density of bound charge on the surface of the sphere.
(b) Define the magnetization of a material and explain what is meant by the term
bound current.
Consider a sample of material with spatially-dependent magnetization M(x) occu-
pying a region V with surface S. Show that, in the absence of free currents, the resulting
vector potential A(x) can be ascribed to bulk and surface densities of bound current.
Consider an infinite cylinder of radius r consisting of a material with permeability µ
surrounded by a region of vacuum. A thin wire carrying current I is placed along the axis
of the cylinder. Determine the direction and magnitude of the resulting bound current
density on the surface of the cylinder. What is the magnetization M(x) on the surface of
the cylinder?

Part II, 2018 List of Questions [TURN OVER

 40

Paper 3, Section II
37D Electrodynamics
Starting from the covariant form of the Maxwell equations and making a suitable
choice of gauge which you should specify, show that the 4-vector potential due to an
arbitrary 4-current J µ (x) obeys the wave equation,
 
1 ∂2
∇ − 2 2 Aµ = −µ0 J µ ,
2
c ∂t

where xµ = (ct, x).

Use the method of Green’s functions to show that, for a localised current distribu-
tion, this equation is solved by
Z
µ µ0 J µ (tret , x′ )
A (t, x) = d3 x′ ,
4π |x − x′ |

for some tret that you should specify.

A point particle, of charge q, moving along a worldline y µ (τ ) parameterised by
proper time τ , produces a 4-vector potential

µ0 qc ẏ µ (τ⋆ )
Aµ (x) =
4π |Rν (τ⋆ )ẏν (τ⋆ )|

where Rµ (τ ) = xµ − y µ (τ ). Define τ⋆ (x) and draw a spacetime diagram to illustrate its

physical significance.
Suppose the particle follows a circular trajectory,

y(t) = (R cos(ω t), R sin(ω t), 0)

(with y 0 = ct), in some inertial frame with coordinates (ct, x, y, z). Evaluate the resulting
4-vector potential at a point on the z-axis as a function of z and t.

Part II, 2018 List of Questions

 41

Paper 2, Section II
38C Fluid Dynamics II
An initially unperturbed two-dimensional inviscid jet in −h < y < h has uniform
speed U in the x direction, while the surrounding fluid is stationary. The unperturbed
velocity field u = (u, v) is therefore given by

u=0 in y > h,
u =U in − h < y < h,
u=0 in y < −h.

Consider separately disturbances in which the layer occupies −h − η < y < h + η (varicose
disturbances) and disturbances in which the layer occupies −h + η < y < h + η (sinuous
disturbances), where η(x, t) = η̂eikx+σt , and determine the dispersion relation σ(k) in each
case.
Find asymptotic expressions for the real part σR of σ in the limits k → 0 and k → ∞
and draw sketches of σR (k) in each case.
Compare the rates of growth of the two types of disturbance.

Paper 1, Section II
38C Fluid Dynamics II
A two-dimensional layer of very viscous fluid of uniform thickness h(t) sits on a
stationary, rigid surface y = 0. It is impacted by a stream of air (which can be assumed
inviscid) such that the air pressure at y = h is p0 − 21 ρa E 2 x2 , where p0 and E are constants,
ρa is the density of the air, and x is the coordinate parallel to the surface.
What boundary conditions apply to the velocity u = (u, v) and stress tensor σ of
the viscous fluid at y = 0 and y = h?
By assuming the form ψ = xf (y) for the stream function of the flow, or otherwise,
solve the Stokes equations for the velocity and pressure fields. Show that the layer thins
at a rate
dh 1 ρa 2 3
V =− = E h .
dt 3 µ

Part II, 2018 List of Questions [TURN OVER

 42

Paper 4, Section II
38C Fluid Dynamics II
A cylinder of radius a rotates about its axis with angular velocity Ω while its axis
is fixed parallel to and at a distance a + h0 from a rigid plane, where h0 ≪ a. Fluid of
kinematic viscosity ν fills the space between the cylinder and the plane. Determine the gap
width h between the cylinder and the plane as a function of a coordinate x parallel to the
surface of the wall and orthogonal to the axis of the cylinder. What is the characteristic
length scale, in the x direction, for changes in the gap width? Taking an appropriate
approximation for h(x), valid in the region where the gap width h is small, use lubrication
theory to determine that the volume flux between the wall and the cylinder (per unit
length along the axis) has magnitude 23 aΩh0 , and state its direction.
Evaluate the tangential shear stress τ on the surface of the cylinder. Approximating
the torque
R ∞ on the cylinder (per unit length along the axis) in the form of an integral
T = a −∞ τ dx, find the torque T to leading order in h0 /a ≪ 1.
3/2
Explain the restriction a1/2 Ωh0 /ν ≪ 1 for the theory to be valid.
Z ∞ Z ∞
dx π dx 3π
[You may use the facts that 2 2
= and 2 3
= .]
−∞ (1 + x ) 2 −∞ (1 + x ) 8

Part II, 2018 List of Questions

 43

Paper 3, Section II
39C Fluid Dynamics II
For two Stokes flows u(1) (x) and u(2) (x) inside the same volume V with different
boundary conditions on its boundary S, prove the reciprocal theorem
Z Z
(1) (2) (2) (1)
ui σij nj dS = ui σij nj dS,
S S

where σ (1) and σ (2) are the stress tensors associated with the flows.
Stating clearly any properties of Stokes flow that you require, use the reciprocal
theorem to prove that the drag F on a body translating with uniform velocity U is given
by
Fi = Aij Uj ,
where A is a symmetric second-rank tensor that depends only on the geometry of the
body.
A slender rod falls slowly through very viscous fluid with its axis inclined to the
vertical. Explain why the rod does not rotate, stating any properties of Stokes flow that
you use.
When the axis of the rod is inclined at an angle θ to the vertical, the centre of mass
of the rod travels at an angle φ to the vertical. Given that the rod falls twice as quickly
when its axis is vertical as when its axis is horizontal, show that
sin θ cos θ
tan φ = .
1 + cos2 θ

Part II, 2018 List of Questions [TURN OVER

 44

Paper 1, Section I
7B Further Complex Methods
The Beta and Gamma functions are defined by
Z 1
B(p, q) = tp−1 (1 − t)q−1 dt,
0
Z ∞
Γ(p) = e−t tp−1 dt,
0

where Re p > 0, Re q > 0.

(a) By using a suitable substitution, or otherwise, prove that

B(z, z) = 21−2z B(z, 12 )

for Re z > 0. Extending B by analytic continuation, for which values of z ∈ C does this
result hold?
(b) Prove that
Γ(p)Γ(q)
B(p, q) = ,
Γ(p + q)
for Re p > 0, Re q > 0.

Paper 2, Section I
7B Further Complex Methods
Show that Z ∞
cos nx − cos mx
dx = π(m − n),
−∞ x2
in the sense of Cauchy principal value, where n and m are positive integers. [State clearly
any standard results involving contour integrals that you use.]

Paper 3, Section I
7B Further Complex Methods
Using a suitable branch cut, show that
Z a
a2 π
(a2 − x2 )1/2 dx = ,
−a 2

where a > 0.

Part II, 2018 List of Questions

 45

Paper 4, Section I
7B Further Complex Methods
State the conditions for a point z = z0 to be a regular singular point of a linear
second-order homogeneous ordinary differential equation in the complex plane.
Find all singular points of the Bessel equation
 
2 ′′ ′ 2 1
z y (z) + zy (z) + z − y(z) = 0 , (∗)
4

and determine whether they are regular or irregular.

√
By writing y(z) = f (z)/ z, find two linearly independent solutions of (∗). Comment
on the relationship of your solutions to the nature of the singular points.

Paper 2, Section II
13B Further Complex Methods
Consider a multi-valued function w(z).
(a) Explain what is meant by a branch point and a branch cut.
(b) Consider z = ew .

(i) By writing z = reiθ , where 0 6 θ < 2π, and w = u + iv, deduce the
expression for w(z) in terms of r and θ. Hence, show that w is infinitely
valued and state its principal value.

(ii) Show that z = 0 and z = ∞ are the branch points of w. Deduce that the
line Im z = 0, Re z > 0 is a possible choice of branch cut.

(iii) Use the Cauchy–Riemann conditions to show that w is analytic in the cut
dw 1
plane. Show that = .
dz z

Part II, 2018 List of Questions [TURN OVER

 46

Paper 1, Section II
14B Further Complex Methods
The equation
zw′′ + 2aw′ + zw = 0, (†)
where a is a constant with Re a > 0, has solutions of the form
Z
w(z) = ezt f (t)dt,
γ

for suitably chosen contours γ and some suitable function f (t).

(a) Find f (t) and determine the condition on γ, which you should express in terms
of z, t and a.
(b) Use the results of part (a) to show that γ can be a finite contour and specify
two possible finite contours with the help of a clearly labelled diagram. Hence, find the
corresponding solution of the equation (†) in the case a = 1.
(c) In the case a = 1 and real z, show that γ can be an infinite contour and specify
two possible infinite contours with the help of a clearly labelled diagram. [Hint: Consider
separately the cases z > 0 and z < 0.] Hence, find a second, linearly independent solution
of the equation (†) in this case.

Part II, 2018 List of Questions

 47

Paper 4, Section II
18I Galois Theory
Let K be a field of characteristic p > 0 and let L be the splitting field of the
polynomial f (t) = tp − t + a over K, where a ∈ K. Let α ∈ L be a root of f (t).
If L 6= K, show that f (t) is irreducible over K, that L = K(α), and that L is a
Galois extension of K. What is Gal(L/K)?

Paper 3, Section II
18I Galois Theory
Let L be a finite field extension of a field K, and let G be a finite group of K-
automorphisms of L. Denote by LG the field of elements of L fixed by the action of G.
(a) Prove that the degree of L over LG is equal to the order of the group G.
(b) For any α ∈ L write f (t, α) = Πg∈G (t − g(α)).

(i) Suppose that L = K(α). Prove that the coefficients of f (t, α) generate LG
over K.

(ii) Suppose that L = K(α1 , α2 ). Prove that the coefficients of f (t, α1 ) and
1/2 1/2
f (t, α2 ) lie in LG . By considering the case L = K(a1 , a2 ) with a1 and
a2 in K, or otherwise, show that they need not generate LG over K.

Paper 2, Section II
18I Galois Theory
Let K be a field and let f (t) be a monic polynomial with coefficients in K. What
is meant by a splitting field L for f (t) over K? Show that such a splitting field exists and
is unique up to isomorphism.
Now suppose that K is a finite field. Prove that L is a Galois extension of K with
cyclic Galois group. Prove also that the degree of L over K is equal to the least common
multiple of the degrees of the irreducible factors of f (t) over K.
Now suppose K is the field with two elements, and let

Pn = {f (t) ∈ K[t] | f has degree n and is irreducible over K}.

How many elements does the set P9 have?

Part II, 2018 List of Questions [TURN OVER

 48

Paper 1, Section II
18I Galois Theory
Let f (t) = t4 + bt2 + ct + d be an irreducible quartic with rational coefficients.
Explain briefly why it is that if the cubic g(t) = t3 + 2bt2 + (b2 − 4d)t − c2 has S3 as its
Galois group then the Galois group of f (t) is S4 .
For which prime numbers p is the Galois group of t4 + pt + p a proper subgroup
of S4 ? [You may assume that the discriminant of t3 + λt + µ is −4λ3 − 27µ2 .]

Part II, 2018 List of Questions

 49

Paper 1, Section II
37E General Relativity
Consider the de Sitter metric

ds2 = −dt2 + e2Ht (dx2 + dy 2 + dz 2 ) ,

where H > 0 is a constant.

(a) Write down the Lagrangian governing the geodesics of this metric. Use the
Euler–Lagrange equations to determine all non-vanishing Christoffel symbols.
(b) Let C be a timelike geodesic parametrized by proper time τ with initial conditions
at τ = 0,
t = 0 , x = y = z = 0 , ẋ = v0 > 0 , ẏ = ż = 0 ,
where the dot denotes differentiation with respect to τ and v0 is a constant. Assuming
both t and τ to be future oriented, show that at τ = 0,
q
ṫ = 1 + v02 .

(c) Find a relation between τ and t along the geodesic of part (b) and show that
t → −∞ for a finite value of τ . [You may use without proof that
Z √
1 1 1 + ae−bu + 1
√ du = ln √ + constant , a, b > 0.]
1 + ae−bu b 1 + ae−bu − 1

(d) Briefly interpret this result.

Part II, 2018 List of Questions [TURN OVER

 50

Paper 2, Section II
37E General Relativity
The Friedmann equations and the conservation of energy-momentum for a spatially
homogeneous and isotropic universe are given by:

ȧ2 + k 2aä + ȧ2 + k ȧ

3 − Λ = 8πρ , − Λ = −8πP , ρ̇ = −3 (P + ρ) ,
a2 a2 a
where a is the scale factor, ρ the energy density, P the pressure, Λ the cosmological
constant and k = +1, 0, −1.
(a) Show that for an equation of state P = wρ, w = constant, the energy density
obeys ρ = 3µ
8π a
−3(1+w) , for some constant µ.

(b) Consider the case of a matter dominated universe, w = 0, with Λ = 0. Write

the equation of motion for the scale factor a in the form of an effective potential equation,

ȧ2 + V (a) = C ,

where you should determine the constant C and the potential V (a). Sketch the potential
V (a) together with the possible values of C and qualitatively discuss the long-term
dynamics of an initially small and expanding universe for the cases k = +1, 0, −1.
(c) Repeat the analysis of part (b), again assuming w = 0, for the cases:

(i) Λ > 0, k = −1,

(ii) Λ < 0, k = 0,

(iii) Λ > 0, k = 1.

Discuss all qualitatively different possibilities for the dynamics of the universe in each case.

Part II, 2018 List of Questions

 51

Paper 4, Section II
37E General Relativity
(a) In the Newtonian weak-field limit, we can write the spacetime metric in the form

ds2 = −(1 + 2Φ)dt2 + (1 − 2Φ)δij dxi dxj , (∗)

where δij dxi dxj = dx2 + dy 2 + dz 2 and the potential Φ(t, x, y, z), as well as the velocity v
of particles moving in the gravitational field are assumed to be small, i.e.,

Φ, ∂t Φ, ∂xi Φ, v 2 ≪ 1.

Use the geodesic equation for this metric to derive the equation of motion for a massive
point particle in the Newtonian limit.
(b) The far-field limit of the Schwarzschild metric is a special case of (∗) given, in
spherical coordinates, by
   
2 2M 2 2M
ds = − 1 − dt + 1 + (dr 2 + r 2 dθ 2 + r 2 sin2 θdϕ2 ) ,
r r

where now M/r ≪ 1. For the following questions, state your results to first order in M/r,
i.e. neglecting terms of O((M/r)2 ).

(i) Let r1 , r2 ≫ M . Calculate the proper length S along the radial curve from r1
to r2 at fixed t, θ, ϕ.

(ii) Consider a massless particle moving radially from r = r1 to r = r2 . According

to an observer at rest at r2 , what time T elapses during this motion?

(iii) The effective velocity of the particle as seen by the observer at r2 is defined as
veff := S/T . Evaluate veff and then take the limit of this result as r1 → r2 .
Briefly discuss the value of veff in this limit.

Part II, 2018 List of Questions [TURN OVER

 52

Paper 3, Section II
38E General Relativity
The Schwarzschild metric in isotropic coordinates x̄ᾱ = (t̄, x̄, ȳ, z̄), ᾱ = 0, . . . , 3, is
given by:

(1 − A)2 2
ds2 = ḡᾱβ̄ dx̄ᾱ dx̄β̄ = − dt̄ + (1 + A)4 (dx̄2 + dȳ 2 + dz̄ 2 )
(1 + A)2

where p
m
A=, r̄ = x̄2 + ȳ 2 + z̄ 2 ,
2r̄
and m is the mass of the black hole.
(a) Let xµ = (t, x, y, z), µ = 0, . . . , 3, denote a coordinate system related to x̄ᾱ by

t̄ = γ(t − vx), x̄ = γ(x − vt), ȳ = y, z̄ = z ,

√
where γ = 1/ 1 − v 2 and −1 < v < 1. Write down the transformation matrix ∂ x̄ᾱ /∂xµ ,
briefly explain its physical meaning and show that the inverse transformation is of the
same form, but with v → −v.
(b) Using the coordinate transformation matrix of part (a), or otherwise, show that
the components gµν of the metric in coordinates xµ are given by

ds2 = gµν dxµ dxν = f (A)(−dt2 + dx2 + dy 2 + dz 2 ) + γ 2 g(A)(dt − v dx)2 ,

where f and g are functions of A that you should determine. You should also express A
in terms of the coordinates (t, x, y, z).
(c) Consider the limit v → 1 with p = mγ held constant. Show that for points x 6= t
the function A → 0, while γ 2 A tends to a finite value, which you should determine. Hence
determine the metric components gµν at points x 6= t in this limit.

Part II, 2018 List of Questions

 53

Paper 4, Section II
17I Graph Theory
Let s > 3. Define the Ramsey number R(s). Show that R(s) exists and that
R(s) 6 4s .
Show that R(3) = 6. Show that (up to relabelling the vertices) there is a unique
way to colour the edges of the complete graph K5 blue and yellow with no monochromatic
triangle.
What is the least positive integer n such that the edges of the complete graph K6
can be coloured blue and yellow in such a way that there are precisely n monochromatic
triangles?

Paper 3, Section II
17I Graph Theory
What does it mean to say that a graph G has a k-colouring? What are the chromatic
number χ(G) and the independence number α(G) of a graph G? For each r > 3, give an
example of a graph G such that χ(G) > r but Kr 6⊂ G.
Let g, k > 3. Show that there exists a graph G containing no cycle of length 6 g
with χ(G) > k.
Show also that if n is sufficiently large then there is a triangle-free G of order n with
α(G) < n0.7 .

Paper 2, Section II
17I Graph Theory
Let G be a graph and A, B ⊂ V (G). Show that if every AB-separator in G has
order at least k then there exist k vertex-disjoint AB-paths in G.
Let k > 3 and assume that G is k-connected. Show that G must contain a cycle of
length at least k.
Assume further that |G| > 2k. Must G contain a cycle of length at least 2k? Justify
your answer.
What is the largest integer n such that any 3-connected graph G with |G| > n must
contain a cycle of length at least n?
[No form of Menger’s theorem or of the max-flow-min-cut theorem may be assumed
without proof.]

Part II, 2018 List of Questions [TURN OVER

 54

Paper 1, Section II
17I Graph Theory
(a) Define ex (n, H) where H is a graph with
at least one edge and n > |H|. Show
that, for any such H, the limit limn→∞ ex (n, H)/ n2 exists.
[You may not assume the Erdős–Stone theorem.]
(b) State the
Erdős–Stone theorem. Use it to deduce that if H is bipartite then
limn→∞ ex (n, H)/ n2 = 0.
 
2− 1t
(c) Let t > 2. Show that ex (n, Kt,t ) = O n .
We say A ⊂ Zn is nice if whenever a, b, c, d ∈ A with a + b = c + d then either
a = c, b = d or a = d, b = c. Let f (n) = max{|A| : A ⊂ Zn , A is nice}. Show that
√
f (n) = O( n).
[Zn denotes the set of integers modulo n, i.e. Zn = {0, 1, 2, . . . , n − 1} with addition
modulo n.]

Part II, 2018 List of Questions

 55

Paper 1, Section II
32A Integrable Systems
Let M = R2n = {(q, p)|q, p ∈ Rn } be equipped with the standard symplectic form
so that the Poisson bracket is given by:
∂f ∂g ∂f ∂g
{f, g} = − ,
∂qj ∂pj ∂pj ∂qj

for f, g real-valued functions on M . Let H = H(q, p) be a Hamiltonian function.

(a) Write down Hamilton’s equations for (M, H), define a first integral of the system
and state what it means that the system is integrable.
(b) State the Arnol’d–Liouville theorem.
(c) Define complex coordinates zj by zj = qj + ipj , and show that if f, g are real-
valued functions on M then:
∂f ∂g ∂g ∂f
{f, g} = −2i + 2i .
∂zj ∂zj ∂zj ∂zj

1
(d) For an n×n anti-Hermitian matrix A with components Ajk , let IA := 2i zj Ajk zk .
Show that:
{IA , IB } = −I[A,B] ,
where [A, B] = AB − BA is the usual matrix commutator.
(e) Consider the Hamiltonian:
1
H = zj zj .
2
Show that (M, H) is integrable and describe the invariant tori.
[In this question j, k = 1, . . . , n, and the summation convention is understood for
these indices.]

Part II, 2018 List of Questions [TURN OVER

 56

Paper 2, Section II
33A Integrable Systems
(a) Let L, A be two families of linear operators, depending on a parameter t, which
act on a Hilbert space H with inner product (, ). Suppose further that for each t, L is
self-adjoint and that A is anti-self-adjoint. State Lax’s equation for the pair L, A, and
show that if it holds then the eigenvalues of L are independent of t.
(b) For ψ, φ : R → C, define the inner product:
Z ∞
(ψ, φ) := ψ(x)φ(x)dx.
−∞

Let L, A be the operators:

 
d3 ψ dψ d
Lψ := i 3 − i q + (qψ) + pψ,
dx dx dx

d2 ψ
Aψ := 3i − 4iqψ,
dx2
where p = p(x, t), q = q(x, t) are smooth, real-valued functions. You may assume that
the normalised eigenfunctions of L are smooth functions of x, t, which decay rapidly as
|x| → ∞ for all t.

(i) Show that if ψ, φ are smooth and rapidly decaying towards infinity then:

(Lψ, φ) = (ψ, Lφ), (Aψ, φ) = −(ψ, Aφ).

Deduce that the eigenvalues of L are real.

(ii) Show that if Lax’s equation holds for L, A, then q must satisfy the Boussinesq
equation:
qtt = aqxxxx + b(q 2 )xx ,
where a, b are constants whose values you should determine. [You may assume
without proof that the identity:
 
dψ d  
LAψ = ALψ − 3i px + (px ψ) + qxxx − 4(q 2 )x ψ,
dx dx

holds for smooth, rapidly decaying ψ.]

Part II, 2018 List of Questions

 57

Paper 3, Section II
33A Integrable Systems
Suppose ψ s : (x, u) 7→ (x̃, ũ) is a smooth one-parameter group of transformations
acting on R2 .
(a) Define the generator of the transformation,

∂ ∂
V = ξ(x, u) + η(x, u) ,
∂x ∂u
where you should specify ξ and η in terms of ψ s .
(b) Define the nth prolongation of V , Pr(n) V and explicitly compute Pr(1) V in
terms of ξ, η.
Recall that if ψ s is a Lie point symmetry of the ordinary differential equation:
 
du dn u
∆ x, u, , . . . , n = 0,
dx dx

then it follows that Pr(n) V [∆] = 0 whenever ∆ = 0.

(c) Consider the ordinary differential equation:

du
= F (x, u) ,
dx
for F a smooth function. Show that if V generates a Lie point symmetry of this equation,
then:
0 = ηx + (ηu − ξx − F ξu ) F − ξFx − ηFu .

(d) Find all the Lie point symmetries of the equation:

du u
= xG 2 ,
dx x
where G is an arbitrary smooth function.

Part II, 2018 List of Questions [TURN OVER

 58

Paper 3, Section II
21F Linear Analysis
(a) Let X be a normed vector space and let Y be a Banach space. Show that the
space of bounded linear operators B(X, Y ) is a Banach space.
(b) Let X and Y be Banach spaces, and let D ⊂ X be a dense linear subspace.
Prove that a bounded linear map T : D → Y can be extended uniquely to a bounded
linear map T : X → Y with the same operator norm. Is the claim also true if one of X
and Y is not complete?
(c) Let X be a normed vector space. Let (xn ) be a sequence in X such that
∞
X
|f (xn )| < ∞ ∀f ∈ X ∗ .
n=1

Prove that there is a constant C such that

∞
X
|f (xn )| 6 Ckf k ∀f ∈ X ∗ .
n=1

Paper 1, Section II
22F Linear Analysis
Let K be a compact Hausdorff space.
(a) State the Arzelà–Ascoli theorem, and state both the real and complex versions
of the Stone–Weierstraß theorem. Give an example of a compact space K and a bounded
set of functions in C(K) that is not relatively compact.
(b) Let f : Rn → R be continuous. Show that there exists a sequence of polynomials
(pi ) in n variables such that

B ⊂ Rn compact ⇒ pi |B → f |B uniformly.

Characterize the set of continuous functions f : Rn → R for which there exists a sequence
of polynomials (pi ) such that pi → f uniformly on Rn .
(c) Prove that if C(K) is equicontinuous then K is finite. Does this implication
remain true if we drop the requirement that K be compact? Justify your answer.

Part II, 2018 List of Questions

 59

Paper 2, Section II
22F Linear Analysis
Let X, Y be Banach spaces and let B(X, Y ) denote the space of bounded linear
operators T : X → Y .
(a) Define what it means for a bounded linear operator T : X → Y to be compact.
Let Ti : X → Y be linear operators with finite rank, i.e., Ti (X) is finite-dimensional.
Assume that the sequence Ti converges to T in B(X, Y ). Show that T is compact.
(b) Let T : X → Y be compact. Show that the dual map T ∗ : Y ∗ → X ∗ is compact.
[Hint: You may use the Arzelà–Ascoli theorem.]
(c) Let X be a Hilbert space and let T : X → X be a compact operator. Let (λj ) be
an infinite sequence of eigenvalues of T with eigenvectors xj . Assume that the eigenvectors
are orthogonal to each other. Show that λj → 0.

Paper 4, Section II
22F Linear Analysis
(a) Let X be a separable normed space. For any sequence (fn )n∈N ⊂ X ∗ with
kfn k 6 1 for all n, show that there is f ∈ X ∗ and a subsequence Λ ⊂ N such that
fn (x) → f (x) for all x ∈ X as n ∈ Λ, n → ∞. [You may use without proof the fact that
X ∗ is complete and that any bounded linear map f : D → R, where D ⊂ X is a dense
linear subspace, can be extended uniquely to an element f ∈ X ∗ .]
(b) Let H be a Hilbert space and U : H → H a unitary map. Let

I = {x ∈ H : U x = x}, W = {U x − x : x ∈ H}.

Prove that I and W are orthogonal, H = I ⊕ W , and that for every x ∈ H,

n−1
1X i
lim U x = P x,
n→∞ n
i=0

where P is the orthogonal projection onto the closed subspace I.

(c) Let T : C(S 1 ) → C(S 1 ) be a linear map, where S 1 = {eiθ ∈ C : θ ∈ R} is the
unit circle, induced by a homeomorphism τ : S 1 → S 1 by (T f )eiθ = f (τ (eiθ )). Prove that
there exists µ ∈ C(S 1 )∗ with µ(1S 1 ) = 1 such that µ(T f ) = µ(f ) for all f ∈ C(S 1 ). (Here
1S 1 denotes the function on S 1 which returns 1 identically.) If T is not the identity map,
does it follow that µ as above is necessarily unique? Justify your answer.

Part II, 2018 List of Questions [TURN OVER

 60

Paper 4, Section II
16G Logic and Set Theory
State and prove the ǫ-Recursion Theorem. [You may assume the Principle of ǫ-
Induction.]
What does it mean to say that a relation r on a set x is well-founded and extensional?
State and prove Mostowski’s Collapsing Theorem. [You may use any recursion theorem
from the course, provided you state it precisely.]
For which sets x is it the case that every well-founded extensional relation on x is
isomorphic to the relation ǫ on some transitive subset of Vω ?

Paper 3, Section II
16G Logic and Set Theory
State and prove the Compactness Theorem for first-order predicate logic. State and
prove the Upward Löwenheim–Skolem Theorem.
[You may assume the Completeness Theorem for first-order predicate logic.]
For each of the following theories, either give axioms (in the specified language) for
the theory or prove that the theory is not axiomatisable.
(i) The theory of finite groups (in the language of groups).
(ii) The theory of groups in which every non-identity element has infinite order (in
the language of groups).
(iii) The theory of total orders (in the language of posets).
(iv) The theory of well-orderings (in the language of posets).
If a theory is axiomatisable by a set S of sentences, and also by a finite set T
of sentences, does it follow that the theory is axiomatisable by some finite subset of S?
Justify your answer.

Paper 2, Section II
16G Logic and Set Theory
State and prove the Knaster–Tarski Fixed-Point Theorem. Deduce the Schröder–
Bernstein Theorem.
Show that the poset P of all countable subsets of R (ordered by inclusion) is not
complete.
Find an order-preserving function f : P → P that does not have a fixed point.
[Hint: Start by well-ordering the reals.]

Part II, 2018 List of Questions

 61

Paper 1, Section II
16G Logic and Set Theory
Give the inductive definition of ordinal exponentiation. Use it to show that αβ 6 αγ
whenever β 6 γ (for α > 1), and also that αβ < αγ whenever β < γ (for α > 2).
Give an example of ordinals α and β with ω < α < β such that αω = β ω .
Show that αβ+γ = αβ αγ , for any ordinals α, β, γ, and give an example to show that
we need not have (αβ)γ = αγ β γ .
For which ordinals α do we have αω1 > ω1 ? And for which do we have αω1 > ω2 ?
Justify your answers.
[You may assume any standard results not concerning ordinal exponentiation.]

Part II, 2018 List of Questions [TURN OVER

 62

Paper 1, Section I
6C Mathematical Biology
Consider a birth-death process in which the birth and death rates in a population
of size n are, respectively, Bn and Dn, where B and D are per capita birth and death
rates.
(a) Write down the master equation for the probability, pn (t), of the population
having size n at time t.
(b) Obtain the differential equations for the rates of change of the mean µ(t) = hni
and the variance σ 2 (t) = hn2 i − hni2 in terms of µ, σ, B and D.
(c) Compare the equations obtained above with the deterministic description of the
evolution of the population size, dn/dt = (B − D)n. Comment on why B and D cannot
be uniquely deduced from the deterministic model but can be deduced from the stochastic
description.

Paper 2, Section I
6C Mathematical Biology
Consider a model of an epidemic consisting of populations of susceptible, S(t), in-
fected, I(t), and recovered, R(t), individuals that obey the following differential equations

dS
= aR − bSI,
dt
dI
= bSI − cI,
dt
dR
= cI − aR,
dt
where a, b and c are constant. Show that the sum of susceptible, infected and recovered
individuals is a constant N . Find the fixed points of the dynamics and deduce the con-
dition for an endemic state with a positive number of infected individuals. Expressing
R in terms of S, I and N , reduce the system of equations to two coupled differential
equations and, hence, deduce the conditions for the fixed point to be a node or a focus.
How do small perturbations of the populations relax to the steady state in each case?

Part II, 2018 List of Questions

 63

Paper 3, Section I
6C Mathematical Biology
Consider a nonlinear model for the axisymmetric dispersal of a population in two
spatial dimensions whose density, n(r, t), obeys

∂n
= D∇ · (n∇n) ,
∂t
where D is a positive constant, r is a radial polar coordinate, and t is time.
Show that Z ∞
2π n(r, t)rdr = N
0
is constant. Interpret this condition.
Show that a similarity solution of the form
 1/2  
N r
n(r, t) = f
Dt (N Dt)1/4

is valid for t > 0 provided that the scaling function f (x) satisfies
 
d df 1 2
xf + x f = 0.
dx dx 4

Show that there exists a value x0 (which need not be evaluated) such that f (x) > 0 for
x < x0 but f (x) = 0 for x > x0 . Determine the area within which n(r, t) > 0 at time t in
terms of x0 .

Hint: The gradient and divergence operators in cylindrical polar coordinates act on
radial functions f and g as
∂f 1 ∂ i
∇f (r) = r̂ , ∇ · [g(r)r̂] = (rg(r)).
∂r r ∂r

Part II, 2018 List of Questions [TURN OVER

 64

Paper 4, Section I
6C Mathematical Biology
Consider a model of a population Nτ in discrete time
rNτ
Nτ +1 = ,
(1 + bNτ )2

where r, b > 0 are constants and τ = 1, 2, 3, . . .. Interpret the constants and show that for
r > 1 there is a stable fixed point.
Suppose the initial condition is N1 = 1/b and that r > 4. Show, using a cobweb
diagram, that the population Nτ is bounded as

4r 2 r
2
6 Nτ 6
(4 + r) b 4b

and attains the bounds.

Part II, 2018 List of Questions

 65

Paper 3, Section II
13C Mathematical Biology
Consider fluctuations of a population described by the vector x = (x1 , x2 , . . . , xN ).
The probability of the state x at time t, P (x, t), obeys the multivariate Fokker–Planck
equation
∂P ∂
1 ∂2

=− Ai (x)P + Bij (x)P ,
∂t ∂xi 2 ∂xi ∂xj
where P = P (x, t), Ai is a drift vector and Bij is a symmetric positive-definite diffusion
matrix, and the summation convention is used throughout.
(a) Show that the Fokker–Planck equation can be expressed as a continuity equation

∂P
+ ∇ · J = 0,
∂t
for some choice of probability flux J which you should determine explicitly. Here,
∇ = ( ∂x∂ 1 , ∂x∂ 2 , . . . , ∂x∂N ) denotes the gradient operator.
(b) Show that the above implies that an initially normalised probability distribution
remains normalised, Z
P (x, t)dV = 1,

at all times, where the volume element dV = dx1 dx2 . . . dxN .

(c) Show that the first two moments of the probability distribution obey

d
hxk i = hAk i
dt
d
hxk xl i = hxl Ak + xk Al + Bkl i.
dt

(d) Now consider small fluctuations with zero mean, and assume that it is possible to
linearise the drift vector and the diffusion matrix as Ai (x) = aij xj and Bij (x) = bij where
aij has real negative eigenvalues and bij is a symmetric positive-definite matrix. Express
the probability flux in terms of the matrices aij and bij and assume that it vanishes in the
stationary state.
(e) Hence show that the multivariate normal distribution,
1 1
P (x) = exp(− Dij xi xj ),
Z 2
where Z is a normalisation and Dij is symmetric, is a solution of the linearised Fokker–
Planck equation in the stationary state, and obtain an equation that relates Dij to the
matrices aij and bij .
(f) Show that the inverse of the matrix Dij is the matrix of covariances Cij = hxi xj i
and obtain an equation relating Cij to the matrices aij and bij .

Part II, 2018 List of Questions [TURN OVER

 66

Paper 4, Section II
14C Mathematical Biology
An activator-inhibitor reaction diffusion system is given, in dimensionless form, by

∂u ∂ 2 u u2 ∂v ∂2v
=d 2 + − 2bu, = + u2 − v,
∂t ∂x v ∂t ∂x2
where d and b are positive constants. Which symbol represents the concentration of
activator and which the inhibitor ? Determine the positive steady states and show, by
an examination of the eigenvalues in a linear stability analysis of the spatially uniform
situation, that the reaction kinetics are stable if b < 21 .
Determine the conditions for the steady state to be driven unstable by diffusion,
and sketch the (b, d) parameter space in which the diffusion-driven instability occurs.
Find the critical wavenumber kc at the bifurcation to such a diffusion-driven instability.

Part II, 2018 List of Questions

 67

Paper 2, Section II
20G Number Fields
Let p ≡ 1 mod 4 be a prime, and let ω = e2πi/p . Let L = Q(ω).
(a) Show that [L : Q] = p − 1.
√
(b) Calculate disc(1, ω, ω 2 , . . . , ω p−2 ). Deduce that p ∈ L.
√
(c) Now suppose p = 5. Prove that OL× = {±ω a ( 12 + 25 )b | a, b ∈ Z}. [You may use
any general result without proof, provided that you state it precisely.]

Paper 4, Section II
20G Number Fields
√
Let m > 2 be a square-free integer, and let n > 2 be an integer. Let L = Q( n m).
(a) By considering the factorisation of (m) into prime ideals, show that [L : Q] = n.
(b) Let T : L × L → Q be the bilinear form defined by T (x, y) = trL/Q (xy). Let
√ ∗
βi = n m i , i = 0, . . . , n − 1. Calculate the dual basis β0∗ , . . . , βn−1 of L with respect to T ,
1 √
and deduce that OL ⊂ nm Z[ m ]. n

√
(c) Show that if p is a prime and n = m = p, then OL = Z[ p p ].

Paper 1, Section II
20G Number Fields
(a) Let m > 2 be an integer such that p = 4m − 1 is prime. Suppose that the ideal
√
class group of L = Q( −p) is trivial. Show that if n > 0 is an integer and n2 +n+m < m2 ,
then n2 + n + m is prime.
√
(b) Show that the ideal class group of Q( −163) is trivial.

Part II, 2018 List of Questions [TURN OVER

 68

Paper 1, Section I
1G Number Theory
(a) State and prove the Chinese remainder theorem.
(b) An integer n is squarefull if whenever p is prime and p|n, then p2 |n. Show that
there exist 1000 consecutive positive integers, none of which are squarefull.

Paper 2, Section I
1G Number Theory
Define the Legendre symbol, and state Gauss’s lemma. Show that if p is an odd
prime, then  
2 2
= (−1)(p −1)/8 .
p
 
105
Use the law of quadratic reciprocity to compute .
149

Paper 3, Section I
1G Number Theory
What is a multiplicative
X function? Show that if f (n) is a multiplicative function,
then so is g(n) = f (d).
d|n

Define the Möbius function µ(n), and show that it is multiplicative. Deduce that
(
X 1 if n = 1
µ(d) =
d|n
0 if n > 1

and that X n

f (n) = µ(e)g .
e
e|n

What is g(n) if f (n) = n? What is f (n) if g(n) = n?

Part II, 2018 List of Questions

 69

Paper 4, Section I
1G Number Theory
Show that if a continued fraction is periodic, then it represents a quadratic irrational.
What number is represented by the continued fraction [7, 7, 7, . . . ]?
√
Compute the continued fraction expansion of 23. Hence or otherwise find a
solution in positive integers to the equation x2 − 23y 2 = 1.

Paper 4, Section II
11G Number Theory
(a) State and prove the Fermat–Euler theorem. Let p be a prime and k a positive
integer. Show that bk ≡ b (mod p) holds for every integer b if and only if k ≡ 1 (mod p−1).
(b) Let N > 3 be an odd integer and b be an integer with (b, N ) = 1. What does it
mean to say that N is a Fermat pseudoprime to base b? What does it mean to say that
N is a Carmichael number ?
Show that every Carmichael number is squarefree, and that if N is squarefree, then
N is a Carmichael number if and only if N ≡ 1 (mod p − 1) for every prime divisor p of
N . Deduce that a Carmichael number is a product of at least three primes.
(c) Let r be a fixed odd prime. Show that there are only finitely many pairs of
primes p, q for which N = pqr is a Carmichael number.
[You may assume throughout that (Z/pn Z)∗ is cyclic for every odd prime p and
every integer n > 1.]

Paper 3, Section II
11G Number Theory
What does it mean to say that a positive definite binary quadratic form is reduced?
What does it mean to say that two binary quadratic forms are equivalent? Show that
every positive definite binary quadratic form is equivalent to some reduced form.
Show that the reduced positive definite binary quadratic forms of discriminant −35
are f1 = x2 +xy +9y 2 and f2 = 3x2 +xy +3y 2 . Show also that a prime p > 7 is represented
by fi if and only if (
p p +1 (i = 1)
= =
5 7 −1 (i = 2).

Part II, 2018 List of Questions [TURN OVER

 70

Paper 4, Section II
40E Numerical Analysis
The inverse discrete Fourier transform Fn−1 : Rn → Rn is given by the formula
n−1
X
x= Fn−1 y, where xℓ = ωnjℓ yj , ℓ = 0, . . . , n−1.
j=0

Here, ωn = exp 2πi p

n is the primitive root of unity of degree n and n = 2 , p = 1, 2, . . ..
−1
(a) Show how to assemble x = F2m y in a small number of operations if the Fourier
transforms of the even and odd parts of y,

x(E) = Fm
−1 (E)
y , x(O) = Fm
−1 (O)
y ,

are already known.

(b) Describe the Fast Fourier Transform (FFT) method for evaluating x, and draw
a diagram to illustrate the method for n = 8.
(c) Find the cost of the FFT method for n = 2p (only multiplications count).
(d) For n = 4 use the FFT method to find x = Fn−1 y when:

(i) y = (1, −1, 1, −1), (ii) y = (1, 1, −1, −1).

Part II, 2018 List of Questions

 71

Paper 2, Section II
40E Numerical Analysis
2
The Poisson equation ddxu2 = f in the unit interval [0, 1], with u(0) = u(1) = 0, is
discretised with the formula

ui−1 − 2ui + ui+1 = h2 fi , 1 6 i 6 n,

where u0 = un+1 = 0, h = (n + 1)−1 , the grid points are at x = ih and ui ≈ u(ih).

(a) Write the above system of equations in the vector form Au = b and describe
the relaxed Jacobi method with relaxation parameter ω for solving this linear system.
(b) For x∗ and x(ν) being the exact and the iterated solution, respectively, let
e(ν) := x(ν) − x∗ be the error and Hω be the iteration matrix, so that

e(ν+1) = Hω e(ν) .

Express Hω in terms of the matrix A and the relaxation parameter ω. Using the fact that
for any n × n Toeplitz symmetric tridiagonal matrix, the eigenvectors vk (k = 1, . . . , n)
have the form:
π
v k = (sin ikx)ni=1 , x = n+1 ,
find the eigenvalues λk (A) of A. Hence deduce the eigenvalues λk (ω) of Hω .
(c) For A as above, let
n
X (ν)
e(ν) = ak v k
k=1

be the expansion of the error with respect to the eigenvectors (v k ) of Hω .

Find the range of the parameter ω which provides convergence of the method for
any n, and prove that, for any such ω, the rate of convergence e(ν) → 0 is not faster than
(1 − c/n2 )ν when n is large.
(ν)
(d) Show that, for an appropriate range of ω, the high frequency components ak
( n+1
2 6 k 6 n) of the error e(ν) tend to zero much faster than the rate obtained in
part (c). Determine the optimal parameter ω∗ which provides the largest supression of the
high frequency components per iteration, and find the corresponding attenuation factor µ∗
(ν+1) (ν)
assuming n is large. That is, find the least µω such that |ak | 6 µω |ak | for n+1
2 6 k 6 n.

Part II, 2018 List of Questions [TURN OVER

 72

Paper 1, Section II
40E Numerical Analysis
(a) Suppose that A is a real n × n matrix, and w ∈ Rn and λ1 ∈ R are given so that
Aw = λ1 w. Further, let S be a non-singular matrix such that Sw = ce(1) , where e(1) is
the first coordinate vector and c 6= 0.
Let Ab = SAS −1 . Prove that the eigenvalues of A are λ1 together with the
b
eigenvalues of the bottom right (n − 1) × (n − 1) submatrix of A.
Explain briefly how, given a vector w, an orthogonal matrix S such that Sw = ce(1)
can be constructed.
(b) Suppose that A is a real n × n matrix, and two linearly independent vectors
v, w ∈ Rn are given such that the linear subspace L{v, w} spanned by v and w is invariant
under the action of A, i.e.,

x ∈ L{v, w} ⇒ Ax ∈ L{v, w}.

Denote by V an n × 2 matrix whose two columns are the vectors v and w, and let S be a
non-singular matrix such that R = SV is upper triangular:
   
v1 w1 r11 r12
 v2 w2   0 r22 
   
SV = S ×   
 v3 w3  =  0 0 
. (∗)
 : :   : : 
vn wn 0 0

b = SAS −1 . Prove that the eigenvalues of A are the eigenvalues of the top left
Again, let A
b together with the eigenvalues of the bottom right (n − 2) × (n − 2)
2 × 2 submatrix of A
b
submatrix of A.
Explain briefly how, for given vectors v, w, an orthogonal matrix S which satisfies
(∗) can be constructed.

Part II, 2018 List of Questions

 73

Paper 3, Section II
41E Numerical Analysis
The diffusion equation for u(x, t):

∂u ∂2u
= , x ∈ R, t > 0,
∂t ∂x2
is solved numerically by the difference scheme

un+1
m = unm + 32 µ(unm−1 − 2unm + unm+1 ) − 12 µ(um−1
n−1 n−1
− 2um n−1
+ um+1 ).
k
Here µ = h2 is the Courant number, with k = ∆t, h = ∆x, and unm ≈ u(mh, nk).
(a) Prove that, as k → 0 with constant µ, the local error of the method is O(k 2 ).
(b) Applying the Fourier stability analysis, show that the method is stable if and
only if µ 6 41 . [Hint: If a polynomial p(x) = x2 − 2αx + β has real roots, then those roots
lie in [a, b] if and only if p(a)p(b) > 0 and α ∈ [a, b].]
(c) Prove that, for the same equation, the leapfrog scheme

un+1
m
n−1
= um + 2µ(unm−1 − 2unm + unm+1 )

is unstable for any choice of µ > 0.

Part II, 2018 List of Questions [TURN OVER

 74

Paper 4, Section II
30K Optimisation and Control
Consider the deterministic system

ẋt = ut

where xt and ut are scalars. Here xt is the state variable and the control variable ut is to
be chosen to minimise, for a fixed h > 0, the cost
Z h
x2h + ct u2t dt,
0

where ct is known and ct > c > 0 for all t. Let F (x, t) be the minimal cost from state x
and time t.
(a) By writing the dynamic programming equation in infinitesimal form and taking
the appropriate limit show that F (x, t) satisfies
 
∂F ∂F
= − inf ct u2 + u , t<h
∂t u ∂x

with boundary condition F (x, h) = x2 .

(b) Determine the form of the optimal control in the special case where ct is constant,
and also in general.

Part II, 2018 List of Questions

 75

Paper 3, Section II
30K Optimisation and Control
The scalars xt , yt , ut are related by the equations

xt = xt−1 + ut−1 , yt = xt−1 + ηt−1 , t = 1, 2, . . . , T,

where the initial state x0 is normally distributed with mean x̂0 and variance 1 and {ηt } is
a sequence of independent random variables each normally distributed with mean 0 and
variance 1. The control variable ut is to be chosen at time t on the basis of information
Wt , where W0 = (x̂0 ) and

Wt = (x̂0 , u0 , . . . , ut−1 , y1 , . . . , yt ), t = 1, 2, . . . , T.

(a) Let x̂1 , x̂2 , . . . , x̂T be the Kalman filter estimates of x1 , x2 , . . . , xT , i.e.

x̂t = x̂t−1 + ut−1 + ht (yt − x̂t−1 )

where ht is chosen to minimise E((x̂t − xt )2 | Wt ). Calculate ht and show that, conditional

on Wt , xt is normally distributed with mean x̂t and variance Vt = 1/(1 + t).
(b) Define


F (WT ) = E x2T | WT , and
 
T
X −1
F (Wt ) = inf E x2T + u2j Wt  , t = 0, . . . , T − 1.
ut ,...,uT −1
j=t

Show that F (Wt ) = x̂2t Pt + dt , where Pt = 1/(T − t + 1), dT = 1/(1 + T ) and

dt−1 = Vt−1 Vt Pt + dt .
(c) Show that the minimising control ut can be expressed in the form ut = −Kt x̂t
and find Kt . How would the expression for Kt be altered if x0 or {ηt } had variances other
than 1?

Part II, 2018 List of Questions [TURN OVER

 76

Paper 2, Section II
30K Optimisation and Control
(a) A ball may be in one of n boxes. A search of the ith box costs ci > 0 and finds
the ball with probability αi > 0 if the ball is in that box. We are given initial probabilities
(Pi1 , i = 1, 2, . . . , n) that the ball is in the ith box.
Show that the policy which at time t = 1, 2, . . . searches the box with the maximal
value of αi Pit /ci minimises the expected searching cost until the ball is found, where Pit is
the probability (given everything that has occurred up to time t) that the ball is in box i.
(b) Next suppose that a reward Ri > 0 is earned if the ball is found in the ith
box. Suppose also that we may decide to stop at any time. Develop the dynamic
programming equation for the value function starting from the probability distribution
(Pit , i = 1, 2, . . . , n).
P
Show that if i ci /(αi Ri ) < 1 then it is never optimal to stop searching until the
ball is found. In this case, is the policy defined in part (a) optimal?

Part II, 2018 List of Questions

 77

Paper 4, Section II
33D Principles of Quantum Mechanics
The spin operators obey the commutation relations [Si , Sj ] = i~ǫijk Sk . Let
|s, σi be an eigenstate of the spin operators Sz and S2 , with Sz |s, σi = σ~ |s, σi and
S2 |s, σi = s(s + 1)~2 |s, σi. Show that
p
S± |s, σi = s(s + 1) − σ(σ ± 1) ~ |s, σ ± 1i ,

where S± = Sx ± iSy . When s = 1, use this to derive the explicit matrix representation
 
0 1 0
~
Sx = √  1 0 1 
2 0 1 0

in a basis in which Sz is diagonal.

A beam of atoms, each with spin 1, is polarised to have spin +~ along the direction
n = (sin θ, 0, cos θ). This beam enters a Stern–Gerlach filter that splits the atoms according
to their spin along the ẑ-axis. Show that N+ /N− = cot4 (θ/2), where N+ (respectively,
N− ) is the number of atoms emerging from the filter with spins parallel (respectively,
anti-parallel) to ẑ.

Part II, 2018 List of Questions [TURN OVER

 78

Paper 1, Section II
33D Principles of Quantum Mechanics
A one-dimensional harmonic oscillator has Hamiltonian
 
† 1
H = ~ω A A +
2
√
where [A, A† ] = 1. Show that A|ni = n|n−1i, where H|ni = (n+ 12 )~ω|ni and hn|ni = 1.
This oscillator is perturbed by adding a new term λX 4 to the Hamiltonian. Given
that
mωX − iP
A= √ ,
2m~ω
show that the ground state of the perturbed system is
r !
~λ √ 3
|0λ i = |0i − 3 2 |2i + |4i ,
4m2 ω 3 2

to first order in λ. [You may use the fact that, in non-degenerate perturbation theory, a
perturbation ∆ causes the first-order shift
X hn|∆|mi
|m(1) i = |ni
Em − En
n6=m

in the mth energy level.]

Part II, 2018 List of Questions

 79

Paper 3, Section II
34D Principles of Quantum Mechanics
A quantum system is prepared in the ground state |0i at time t = 0. It is subjected
to a time-varying Hamiltonian H = H0 + ∆(t). Show that, to first order in ∆(t), the
system evolves as X
|ψ(t)i = ck (t) e−iEk t/~ |ki ,
k

where H0 |ki = Ek |ki and

Z t
1 ′
ck (t) = hk|∆(t′ )|0i ei(Ek −E0 )t /~ dt′ .
i~ 0

A large number of hydrogen atoms, each in the ground state, are subjected to an
electric field (
0 for t < 0
E(t) =
ẑ E0 exp(−t/τ ) for t > 0 ,
where E0 is a constant. Show that the fraction of atoms found in the state |n, ℓ, mi =
|2, 1, 0i is, after a long time and to lowest non-trivial order in E0 ,

215 a20 e2 E02

,
310 ~2 (ω 2 + 1/τ 2 )

where ~ω is the energy difference between the |2, 1, 0i and |1, 0, 0i states, and e is the
electron charge and a0 the Bohr radius. What fraction of atoms lie in the |2, 0, 0i state?
[Hint: You may assume the hydrogenic wavefunctions
   
2 1 r 1 1 r r
hr|1, 0, 0i = √ exp − and hr|2, 1, 0i = √ cos θ exp −
4π a3/2
0
a0 4π (2a0 )3/2 a0 2a0

and the integral Z ∞

m!
r m e−αr dr =
0 αm+1
for m a positive integer.]

Part II, 2018 List of Questions [TURN OVER

 80

Paper 2, Section II
34D Principles of Quantum Mechanics
Explain what is meant by the intrinsic parity of a particle.
In each of the decay processes below, parity is conserved.
A deuteron (d+ ) has intrinsic parity ηd = +1 and spin s = 1. A negatively charged
pion (π − ) has spin s = 0. The ground state of a hydrogenic ‘atom’ formed from a deuteron
and a pion decays to two identical neutrons (n), each of spin s = 21 and parity ηn = +1.
Deduce the intrinsic parity of the pion.
3
The ∆− particle has spin s = 2 and decays as

∆− → π − + n .

What are the allowed values of the orbital angular momentum? In the centre of mass
frame, the vector rπ − rn joining the pion to the neutron makes an angle θ to the ẑ-axis.
The final state is an eigenstate of Jz and the spatial probability distribution is proportional
to cos2 θ. Deduce the intrinsic parity of the ∆− .
[Hint: You may use the fact that the first three Legendre polynomials are given by
1
P0 (x) = 1 , P1 (x) = x , P2 (x) = (3x2 − 1) . ]
2

Part II, 2018 List of Questions

 81

Paper 4, Section II
28K Principles of Statistics
Let g : R → R be an unknown function, twice continuously differentiable with
′′
|g (x)| 6 M for all x ∈ R. For some x0 ∈ R, we know the value g(x0 ) and we wish
to estimate its derivative g ′ (x0 ). To do so, we have access to a pseudo-random number
∗ i.i.d. uniform over [0, 1], and a machine that takes input
generator that gives U1∗ , . . . , UN
x1 , . . . , xN ∈ R and returns g(xi ) + εi , where the εi are i.i.d. N (0, σ 2 ).
(a) Explain how this setup allows us to generate N independent Xi = x0 + hZi ,
where the Zi take value 1 or −1 with probability 1/2, for any h > 0.
(b) We denote by Yi the output g(Xi ) + εi . Show that for some independent ξi ∈ R

h2 ′′
Yi − g(x0 ) = hZi g′ (x0 ) + g (ξi ) + εi .
2

(c) Using the intuition given by the least-squares estimator, justify the use of the
estimator ĝN given by
N
1 X Zi (Yi − g(x0 ))
ĝN = .
N h
i=1

(d) Show that

h2 M 2 σ2
E[|ĝN − g′ (x0 )|2 ] 6 + .
4 N h2
Show that for some choice hN of parameter h, this implies
σM
E[|ĝN − g ′ (x0 )|2 ] 6 √ .
N

Part II, 2018 List of Questions [TURN OVER

 82

Paper 3, Section II
28K Principles of Statistics
In the model {N (θ, Ip ), θ ∈ Rp } of a Gaussian distribution in dimension p, with
unknown mean θ and known identity covariance matrix Ip , we estimate θ based on a
sample of i.i.d. observations X1 , . . . , Xn drawn from N (θ0 , Ip ).
(a) Define the Fisher information I(θ0 ), and compute it in this model.
(b) We recall that the observed Fisher information in (θ) is given by
n
1X
in (θ) = ∇θ log f (Xi , θ) ∇θ log f (Xi , θ)⊤ .
n
i=1

Find the limit of în = in (θ̂M LE ), where θ̂M LE is the maximum likelihood estimator of θ
in this model.
(c) Define the Wald statistic Wn (θ) and compute it. Give the limiting distribution
of Wn (θ0 ) and explain how it can be used to design a confidence interval for θ0 .
[You may use results from the course provided that you state them clearly.]

Paper 2, Section II
28K Principles of Statistics
We consider the model {N (θ, Ip ), θ ∈ Rp } of a Gaussian distribution in dimension
p > 3, with unknown mean θ and known identity covariance matrix Ip . We estimate θ
based on one observation X ∼ N (θ, Ip ), under the loss function

ℓ(θ, δ) = kθ − δk22 .

(a) Define the risk of an estimator θ̂. Compute the maximum likelihood estimator
θ̂M LE of θ and its risk for any θ ∈ Rp .
(b) Define what an admissible estimator is. Is θ̂M LE admissible?
(c) For any c > 0, let πc (θ) be the prior N (0, c2 Ip ). Find a Bayes optimal estimator
θ̂c under this prior with the quadratic loss, and compute its Bayes risk.
(d) Show that θ̂M LE is minimax.
[You may use results from the course provided that you state them clearly.]

Part II, 2018 List of Questions

 83

Paper 1, Section II
29K Principles of Statistics
A scientist wishes to estimate the proportion θ ∈ (0, 1) of presence of a gene in
a population of flies of size n. Every fly receives a chromosome from each of its two
parents, each carrying the gene A with probability (1 − θ) or the gene B with probability
θ, independently. The scientist can observe if each fly has two copies of the gene A (denoted
by AA), two copies of the gene B (denoted by BB) or one of each (denoted by AB). We
let nAA , nBB , and nAB denote the number of each observation among the n flies.
(a) Give the probability of each observation as a function of θ, denoted by f (X, θ),
for all three values X = AA, BB, or AB.
(b) For a vector w = (wAA , wBB , wAB ), we let θ̂w denote the estimator defined by
nAA nBB nAB
θ̂w = wAA + wBB + wAB .
n n n

Find the unique vector w∗ such that θ̂w∗ is unbiased. Show that θ̂w∗ is a consistent
estimator of θ.
(c) Compute the maximum likelihood estimator of θ in this model, denoted by
√
θ̂M LE . Find the limiting distribution of n(θ̂M LE − θ). [You may use results from the
course, provided that you state them clearly.]

Part II, 2018 List of Questions [TURN OVER

 84

Paper 4, Section II
26J Probability and Measure
Let (X, A) be a measurable space. Let T : X → X be a measurable map, and µ a
probability measure on (X, A).
(a) State the definition of the following properties of the system (X, A, µ, T ):

(i) µ is T -invariant.

(ii) T is ergodic with respect to µ.

(b) State the pointwise ergodic theorem.

(c) Give an example of a probability measure preserving system (X, A, µ, T ) in which
Card(T −1 {x}) > 1 for µ-a.e. x.
(d) Assume X is finite and A is the boolean algebra of all subsets of X. Suppose
that µ is a T -invariant probability measure on X such that µ({x}) > 0 for all x ∈ X.
Show that T is a bijection.
(e) Let X = N, the set of positive integers, and A be the σ-algebra of all subsets of
X. Suppose that µ is a T -invariant ergodic probability measure on X. Show that there is
a finite subset Y ⊆ X with µ(Y ) = 1.

Part II, 2018 List of Questions

 85

Paper 2, Section II
26J Probability and Measure
Let (Ω, F, P) be a probability space. Let (Xn )n>1 be a sequence of random variables
with E(|Xn |2 ) 6 1 for all n > 1.
(a) Suppose Z is another random variable such that E(|Z|2 ) < ∞. Why is ZXn
integrable for each n?
(b) Assume E(ZXn ) −−−→ 0 for every random variable Z on (Ω, F, P) such that
n→∞
E(|Z|2 ) < ∞. Show that there is a subsequence Yk := Xnk , k > 1, such that
N
1 X
Yk −−−−→ 0 in L2 .
N N →∞
k=1

(c) Assume that Xn → X in probability. Show that X ∈ L2 . Show that Xn → X

in L1 . Must it converge also in L2 ? Justify your answer.
(d) Assume that the (Xn )n>1 are independent. Give a necessary and sufficient
condition on the sequence (E(Xn )n>1 ) for the sequence
N
1 X
YN = Xk
N
k=1

to converge in L2 .

Part II, 2018 List of Questions [TURN OVER

 86

Paper 3, Section II
26J Probability and Measure
Let m be the Lebesgue measure on the real line. Recall that if E ⊆ R is a Borel
subset, then ( )
X [
m(E) = inf |In |, E ⊆ In ,
n>1 n>1

where the infimum is taken over all covers of E by countably many intervals, and |I|
denotes the length of an interval I.
(a) State the definition of a Borel subset of R.
(b) State a definition of a Lebesgue measurable subset of R.
(c) Explain why the following sets are Borel and compute their Lebesgue measure:
\h1 i
Q, R \ Q, ,n .
n
n>2

(d) State the definition of a Borel measurable function f : R → R.

(e) Let f be a Borel measurable function f : R → R. Is it true that the subset of
all x ∈ R where f is continuous at x is a Borel subset? Justify your answer.
(f) Let E ⊆ [0, 1] be a Borel subset with m(E) = 1/2 + α, α > 0. Show that

E − E := {x − y : x, y ∈ E}

contains the interval (−2α, 2α).

(g) Let E ⊆ R be a Borel subset such that m(E) > 0. Show that for every ε > 0,
there exists a < b in R such that

m(E ∩ (a, b)) > (1 − ε)m((a, b)).

Deduce that E−E contains an open interval around 0.

Part II, 2018 List of Questions

 87

Paper 1, Section II
27J Probability and Measure
(a) Let X be a real random variable with E(X 2 ) < ∞. Show that the variance of
X is equal to inf (E(X − a)2 ).
a∈R

(b) Let f (x) be the indicator function of the interval [−1, 1] on the real line.
Compute the Fourier transform of f .
(c) Show that
Z +∞  2
sin x π
dx = .
0 x 2

(d) Let X be a real random variable and µc

X be its characteristic function.

(i) Assume that |µcX (u)| = 1 for some u ∈ R. Show that there exists θ ∈ R
such that almost surely:

uX ∈ θ + 2πZ.

(ii) Assume that |µc

X (u)| = |µcX (v)| = 1 for some real numbers u, v not equal to
0 and such that u/v is irrational. Prove that X is almost surely constant.
[Hint: You may wish to consider an independent copy of X.]

Part II, 2018 List of Questions [TURN OVER

 88

Paper 4, Section I
10D Quantum Information and Computation
Let Bn denote the set of all n-bit strings. Suppose we are given a 2-qubit quantum
gate Ix0 which is promised to be of the form

| xi x 6= x0
Ix0 | xi =
− | xi x = x0

but the 2-bit string x0 is unknown to us. We wish to determine x0 with the least
number of queries to Ix0 . Define A = I − 2 | ψihψ|, where I is the identity operator
1P
and | ψi = 2 x∈B2 | xi.
(a) Is A unitary? Justify your answer.
(b) Compute the action of Ix0 on | ψi, and the action of | ψihψ| on | x0 i, in each case
expressing your answer in terms of | ψi and | x0 i. Hence or otherwise show that x0 may
be determined with certainty using only one application of the gate Ix0 , together with any
other gates that are independent of x0 .
(c) Let fx0 : B2 → B1 be the function having value 0 for all x 6= x0 and having value
1 for x = x0 . It is known that a single use of Ix0 can be implemented with a single query to
a quantum oracle for the function fx0 . But suppose instead that we have a classical oracle
for fx0 , i.e. a black box which, on input x, outputs the value of fx0 (x). Can we determine
x0 with certainty using a single query to the classical oracle? Justify your answer.

Part II, 2018 List of Questions

 89

Paper 3, Section I
10D Quantum Information and Computation
Let Bn denote the set of all n-bit strings. For any Boolean function on 2 bits
f : B2 → B1 consider the linear operation on 3 qubits defined by

Uf | xi | yi = | xi | y ⊕ f (x)i

for all x ∈ B2 , y ∈ B1 and ⊕ denoting addition of bits modulo 2. Here the first register
is a 2-qubit register and the second is a 1-qubit register. We are able to apply only the
1-qubit Pauli X and Hadamard H gates to any desired qubits, as well as the 3-qubit gate
Uf to any three qubits. We can also perform measurements in the computational basis.
(a) Describe how we can construct the state
1 X
| fi = (−1)f (x) | xi
2
x∈B2

starting from the standard 3-qubit state | 0i | 0i | 0i.

(b) Suppose now that the gate Uf is given to us but f is not specified. However f
is promised to be one of two following cases:

(i) f is a constant function (i.e. f (x) = 0 for all x, or f (x) = 1 for all x),

(ii) for any 2-bit string x = b1 b2 we have f (b1 b2 ) = b1 ⊕ b2 (with ⊕ as above).

Show how we may determine with certainty which of the two cases (i) or (ii) applies, using
only a single application of Uf .

Paper 2, Section I
10D Quantum Information and Computation
(a) The classical controlled-N OT operation applied to the 2-bit string b0 (for b = 0
or 1) achieves the cloning of b, i.e. the result is bb. Let CX denote the quantum controlled-
X (or controlled-N OT ) operation on two qubits. For which qubit states | ψi = a | 0i+b | 1i
will the application of CX to | ψi | 0i (with the first qubit being the control qubit) achieve
the cloning of | ψi? Justify your answer.
(b) Let | α0 i and | α1 i be two distinct non-orthogonal quantum states. State and
prove the quantum no-cloning theorem for unitary processes.

Part II, 2018 List of Questions [TURN OVER

 90

Paper 1, Section I
10D Quantum Information and Computation
(a) Define what it means for a 2-qubit state | ψiAB of a composite quantum
system AB to be entangled.
Consider the 2-qubit state
1  
| αi = √ 2 | 00i − H ⊗ H | 11i
3
where H is the Hadamard gate. From the definition of entanglement, show that | αi is an
entangled state.
(b) Alice and Bob are distantly separated in space. Initially they each hold one
qubit of the 2-qubit entangled state
1  
φ+ = √ | 00i + | 11i .
2
They are able to perform local quantum operations (unitary gates and measurements) on
quantum systems they hold. Alice wants to communicate two classical bits of information
to Bob. Explain how she can achieve this (within their restricted operational resources)
by sending him a single qubit.

Part II, 2018 List of Questions

 91

Paper 2, Section II
15D Quantum Information and Computation
(a) Suppose that Alice and Bob are distantly separated in space and each has one
qubit of the 2-qubit state | φ+ i = √12 (| 00i + | 11i). They also have the ability to perform
local unitary quantum operations and local computational basis measurements, and to
communicate only classically. Alice has a 1-qubit state | αi (whose identity is unknown to
her) which she wants to communicate to Bob. Show how this can be achieved using only
the operational resources, listed above, that they have available.
Suppose now that a third party, called Charlie, joins Alice and Bob. They are all
mutually distantly separated in space and each holds one qubit of the 3-qubit state
1

| γi = √ | 000i + | 111i .
2
As previously with Alice and Bob, they are able to communicate with each other only
classically, e.g. by telephone, and they can each also perform only local unitary operations
and local computational basis measurements. Alice and Bob phone Charlie to say that
they want to do some quantum teleportation and they need a shared | φ+ i state (as defined
above). Show how Charlie can grant them their wish (with certainty), given their joint
possession of | γi and using only their allowed operational resources. [Hint: It may be
useful to consider application of an appropriate Hadamard gate action.]
(b) State the quantum no-signalling principle for a bipartite state | ψiAB of the
composite system AB.
Suppose we are given an unknown one of the two states


| φ+ iAB = √12 | 00iAB + | 11iAB ,


| φ− iAB = √12 | 00iAB − | 11iAB ,

and we wish to identify which state we have. Show that the minimum error probability
for this state discrimination task is zero.
Suppose now that we have access only to qubit B of the received state. Show that
we can now do no better in the state discrimination task than just making a random guess
as to which state we have.

Part II, 2018 List of Questions [TURN OVER

 92

Paper 3, Section II
15D Quantum Information and Computation
In this question you may assume the following fact about the quantum Fourier
transform QF T mod N : if N = Ar and 0 6 x0 < r, where A, r, x0 ∈ Z, then
A−1 r−1
1 X 1 X x0 lA
QF T √ | x0 + kri = √ ω | lAi
A k=0 r
l=0

where ω = e2πi/N .
(a) Let ZN denote the integers modulo N . Let f : ZN → Z be a periodic function
with period r and with the property that f is one-to-one within each period. We have one
instance of the quantum state
N −1
1 X
| fi = √ | xi | f (x)i
N x=0

and the ability to calculate the function f on at most two x values of our choice.
Describe a procedure that may be used to determine the period r with success
probability O(1/ log log N ). As a further requirement, at the end of the procedure we
should know if it has been successful, or not, in outputting the correct period value.
[You may assume that the number of integers less than N that are coprime to N is
O(N/ log log N )].
(b) Consider the function f : Z12 → Z10 defined by f (x) = 3x mod 10.

(i) Show that f is periodic and find the period.

P
(ii) Suppose we are given the state | f i = √112 11
x=0 | xi | f (x)i and we measure
the second register. What are the possible resulting measurement values y
and their probabilities?

(iii) Suppose the measurement result was y = 3. Find the resulting state | αi
of the first register after the measurement.

(iv) Suppose we measure the state QF T | αi (with | αi from part (iii)). What
is the probability of each outcome 0 6 c 6 11?

Part II, 2018 List of Questions

 93

Paper 1, Section II
19I Representation Theory
(a) Define the derived subgroup, G′ , of a finite group G. Show that if χ is a linear
character of G, then G′ 6 ker χ. Prove that the linear characters of G are precisely the
lifts to G of the irreducible characters of G/G′ . [You should state clearly any additional
results that you require.]
(b) For n > 1, you may take as given that the group

G6n := ha, b : a2n = b3 = 1, a−1 ba = b−1 i

has order 6n.

(i) Let ω = e2πi/3 . Show that if ε is any (2n)-th root of unity in C, then there
is a representation of G6n over C which sends
   
0 ε ω 0
a 7→ , b 7→ .
ε 0 0 ω2

(ii) Find all the irreducible representations of G6n .

(iii) Find the character table of G6n .

Part II, 2018 List of Questions [TURN OVER

 94

Paper 2, Section II
19I Representation Theory
(a) Suppose H is a subgroup of a finite group G, χ is an irreducible character
of G and ϕ1 , . . . , ϕr are the irreducible characters of H. Show that in the restriction
χ ↓H = a1 ϕ1 + · · · + ar ϕr , the multiplicities a1 , . . . , ar satisfy
r
X
a2i 6 |G : H|. (†)
i=1

Determine necessary and sufficient conditions under which the inequality in (†) is actually
an equality.
(b) Henceforth suppose that H is a (normal) subgroup of index 2 in G, and that χ
is an irreducible character of G.
Lift the non-trivial linear character of G/H to obtain a linear character of G which
satisfies 
1 if g ∈ H
λ(g) = .
−1 if g 6∈ H

(i) Show that the following are equivalent:

(1) χ ↓H is irreducible;
(2) χ(g) 6= 0 for some g ∈ G with g 6∈ H;
(3) the characters χ and χλ of G are not equal.

(ii) Suppose now that χ ↓H is irreducible. Show that if ψ is an irreducible char-

acter of G which satisfies

ψ ↓H = χ ↓H ,

then either ψ = χ or ψ = χλ.

(iii) Suppose that χ ↓H is the sum of two irreducible characters of H, say

χ ↓H = ψ1 + ψ2 . If φ is an irreducible character of G such that φ ↓H has
ψ1 or ψ2 as a constituent, show that φ = χ.

(c) Suppose that G is a finite group with a subgroup K of index 3, and let χ be an
irreducible character of G. Prove that

hχ ↓K , χ ↓K iK = 1, 2 or 3.

Give examples to show that each possibility can occur, giving brief justification in each
case.

Part II, 2018 List of Questions

 95

Paper 3, Section II
19I Representation Theory
State the row orthogonality relations. Prove that if χ is an irreducible character of
the finite group G, then χ(1) divides the order of G.
Stating clearly any additional results you use, deduce the following statements:
(i) Groups of order p2 , where p is prime, are abelian.
(ii) If G is a group of order 2p, where p is prime, then either the degrees of the
irreducible characters of G are all 1, or they are

1, 1, 2, . . . , 2 (with (p − 1)/2 of degree 2).

(iii) No simple group has an irreducible character of degree 2.

(iv) Let p and q be prime numbers with p > q, and let G be a non-abelian group of
order pq. Then q divides p − 1 and G has q + ((p − 1)/q) conjugacy classes.

Paper 4, Section II
19I Representation Theory
Define G =SU(2) and write down a complete list

{Vn : n = 0, 1, 2, . . .}

of its continuous finite-dimensional irreducible representations. You should define all the
terms you use but proofs are not required. Find the character χVn of Vn . State the
Clebsch–Gordan formula.
(a) Stating clearly any properties of symmetric powers that you need, decompose
the following spaces into irreducible representations of G:

(i) V4 ⊗ V3 , V3 ⊗ V3 , S 2 V3 ;

(ii) V1 ⊗ · · · ⊗ V1 (with n multiplicands);

(iii) S 3 V2 .

(b) Let G act on the space M3 (C) of 3 × 3 complex matrices by

A : X 7→ A1 XA−1
1 ,
 
A 0
where A1 is the block matrix . Show that this gives a representation of G and
0 1
decompose it into irreducible summands.

Part II, 2018 List of Questions [TURN OVER

 96

Paper 2, Section II
23F Riemann Surfaces
State the uniformisation theorem. List without proof the Riemann surfaces which
are uniformised by C∞ and those uniformised by C.
Let U be a domain in C whose complement consists of more than one point. Deduce
that U is uniformised by the open unit disk.
Let R be a compact Riemann surface of genus g and P1 , . . . , Pn be distinct points
of R. Show that R \ {P1 , . . . , Pn } is uniformised by the open unit disk if and only if
2g − 2 + n > 0, and by C if and only if 2g − 2 + n = 0 or −1.
Let Λ be a lattice and X = C/Λ a complex torus. Show that an analytic map
f : C → X is either surjective or constant.
Give with proof an example of a pair of Riemann surfaces which are homeomorphic
but not conformally equivalent.

Paper 3, Section II
23F Riemann Surfaces
Define the degree of an analytic map of compact Riemann surfaces, and state the
Riemann–Hurwitz formula.
Let Λ be a lattice in C and E = C/Λ the associated complex torus. Show that the
map
ψ : z + Λ 7→ −z + Λ
is biholomorphic with four fixed points in E.
Let S = E/ ∼ be the quotient surface (the topological surface obtained by identi-
fying z + Λ and ψ(z + Λ) ), and let p : E → S be the associated projection map. Denote
by E ′ the complement of the four fixed points of ψ, and let S ′ = p(E ′ ). Describe briefly a
family of charts making S ′ into a Riemann surface, so that p : E ′ → S ′ is a holomorphic
map.
Now assume that, by adding finitely many points, it is possible to compactify S ′ to
a Riemann surface S so that p extends to a regular map E → S. Find the genus of S.

Part II, 2018 List of Questions

 97

Paper 1, Section II
24F Riemann Surfaces
Given a complete analytic function F on a domain G ⊂ C, define the germ of a
function element (f, D) of F at z ∈ D. Let G be the set of all germs of function elements
in G. Describe without proofs the topology and complex structure on G and the natural
covering map π : G → G. Prove that the evaluation map E : G → C defined by

E([f ]z ) = f (z)

is analytic on each component of G.

Suppose f : R → S is an analytic map of compact Riemann surfaces with B ⊂ S
the set of branch points. Show that f : R \ f −1 (B) → S \ B is a regular covering map.
Given P ∈ S \ B, explain how any closed curve in S \ B with initial and final points
P yields a permutation of the set f −1 (P ). Show that the group H obtained from all such
closed curves is a transitive subgroup of the group of permutations of f −1 (P ).
Find the group H for the analytic map f : C∞ → C∞ where f (z) = z 2 + z −2 .

Part II, 2018 List of Questions [TURN OVER

 98

Paper 4, Section I

Part II, 2018 List of Questions

 99

5J Statistical Modelling
A scientist is studying the effects of a drug on the weight of mice. Forty mice are
divided into two groups, control and treatment. The mice in the treatment group are
given the drug, and those in the control group are given water instead. The mice are kept
in 8 different cages. The weight of each mouse is monitored for 10 days, and the results
of the experiment are recorded in the data frame Weight.data. Consider the following R
code and its output.

> head(Weight.data)
Time Group Cage Mouse Weight
1 1 Control 1 1 24.77578
2 2 Control 1 1 24.68766
3 3 Control 1 1 24.79008
4 4 Control 1 1 24.77005
5 5 Control 1 1 24.65092
6 6 Control 1 1 24.38436
> mod1 = lm(Weight ~ Time*Group + Cage, data=Weight.data)
> summary(mod1)

Call:
lm(formula = Weight ~ Time * Group + Cage, data = Weight.data)

Residuals:
Min 1Q Median 3Q Max
-1.36903 -0.33527 -0.01719 0.38807 1.24368

Coefficients:
Estimate Std. Error t value Pr(>|t|)
(Intercept) 24.534771 0.100336 244.525 < 2e-16 ***
Time -0.006023 0.012616 -0.477 0.63334
GroupTreatment 0.321837 0.121993 2.638 0.00867 **
Cage2 -0.400228 0.095875 -4.174 3.68e-05 ***
Cage3 0.286941 0.102494 2.800 0.00537 **
Cage4 0.007535 0.095875 0.079 0.93740
Cage6 0.124767 0.125530 0.994 0.32087
Cage8 -0.295168 0.125530 -2.351 0.01920 *
Time:GroupTreatment -0.173515 0.017842 -9.725 < 2e-16 ***
---
Signif. codes: 0 ‘***’ 0.001 ‘**’ 0.01 ‘*’ 0.05 ‘.’ 0.1 ‘ ’ 1

Residual standard error: 0.5125 on 391 degrees of freedom

Multiple R-squared: 0.5591,Adjusted R-squared: 0.55
F-statistic: 61.97 on 8 and 391 DF, p-value: < 2.2e-16

Which parameters describe the rate of weight loss with time in each group?
According to the R output, is there a statistically significant weight loss with time in

Part II, 2018 List of Questions [TURN OVER

 100

the control group?

Three diagnostic plots were generated using the following R code.

mouse1 = (Weight.data$Mouse==1)
plot(Weight.data$Time[mouse1],mod1$residuals[mouse1])
mouse2 = (Weight.data$Mouse==2)
plot(Weight.data$Time[mouse2],mod1$residuals[mouse2])
mouse3 = (Weight.data$Mouse==3)
plot(Weight.data$Time[mouse3],mod1$residuals[mouse3])
mod1$residuals[mouse1]

0.2
0.1
0.0
−0.1

2 4 6 8 10

Weight.data$Time[mouse1]
0.5
mod1$residuals[mouse2]

0.3
0.1
−0.1

2 4 6 8 10

Weight.data$Time[mouse2]

Part II, 2018 List of Questions

 101

mod1$residuals[mouse3]

−0.5
−0.6
−0.7
−0.8

2 4 6 8 10

Weight.data$Time[mouse3]

Based on these plots, should you trust the significance tests shown in the output of
the command summary(mod1)? Explain.

Part II, 2018 List of Questions [TURN OVER

 102

Paper 3, Section I
5J Statistical Modelling
The data frame Cases.of.flu contains a list of cases of flu recorded in 3 London
hospitals during each month of 2017. Consider the following R code and its output.

> table(Cases.of.flu)
Hospital
Month A B C
April 10 40 27
August 9 34 19
December 24 129 81
February 49 134 74
January 45 138 78
July 5 47 35
June 11 36 22
March 20 82 41
May 5 43 23
November 17 82 62
October 6 26 19
September 6 40 21
> Cases.of.flu.table = as.data.frame(table(Cases.of.flu))
> head(Cases.of.flu.table)
Month Hospital Freq
1 April A 10
2 August A 9
3 December A 24
4 February A 49
5 January A 45
6 July A 5
> mod1 = glm(Freq ~., data=Cases.of.flu.table, family=poisson)
> mod1$dev
[1] 28.51836
> levels(Cases.of.flu$Month)
[1] "April" "August" "December" "February" "January" "July"
[7] "June" "March" "May" "November" "October" "September"
> levels(Cases.of.flu$Month) <- c("Q2","Q3","Q4","Q1","Q1","Q3",
+ "Q2","Q1","Q2","Q4","Q4","Q3")
> Cases.of.flu.table = as.data.frame(table(Cases.of.flu))
> mod2 = glm(Freq ~., data=Cases.of.flu.table, family=poisson)
> mod2$dev
[1] 17.9181

Describe a test for the null hypothesis of independence between the variables Month
and Hospital using the deviance statistic. State the assumptions of the test.
Perform the test at the 1% level for each of the two different models shown above.
You may use the table below showing 99th percentiles of the χ2p distribution with a range of
degrees of freedom p. How would you explain the discrepancy between their conclusions?

Part II, 2018 List of Questions

 103

Degrees of freedom 99th percentile Degrees of freedom 99th percentile

1 6.63 21 38.93
2 9.21 22 40.29
3 11.34 23 41.64
4 13.28 24 42.98
5 15.09 25 44.31
6 16.81 26 45.64
7 18.48 27 46.96
8 20.09 28 48.28
9 21.67 29 49.59
10 23.21 30 50.89
11 24.72 31 52.19
12 26.22 32 53.49
13 27.69 33 54.78
14 29.14 34 56.06
15 30.58 35 57.34
16 32.00 36 58.62
17 33.41 37 59.89
18 34.81 38 61.16
19 36.19 39 62.43
20 37.57 40 63.69

Paper 2, Section I
5J Statistical Modelling
Consider a linear model Y = Xβ + σ 2 ε with ε ∼ N (0, I), where the design matrix
X is n by p. Provide an expression for the F -statistic used to test the hypothesis
βp0 +1 = βp0 +2 = · · · = βp = 0 for p0 < p. Show that it is a monotone function of a
log-likelihood ratio statistic.

Part II, 2018 List of Questions [TURN OVER

 104

Paper 1, Section I
5J Statistical Modelling
The data frame Ambulance contains data on the number of ambulance requests
from a Cambridgeshire hospital on different days. In addition to the number of ambulance
requests on each day, the dataset records whether each day fell in the winter season, on a
weekend, or on a bank holiday, as well as the pollution level on each day.

> head(Ambulance)
Winter Weekend Bank.holiday Pollution.level Ambulance.requests
1 Yes Yes No High 16
2 No Yes No Low 7
3 No No No High 22
4 No Yes No Medium 11
5 Yes Yes No High 18
6 No No No Medium 25

A health researcher fitted two models to the dataset above using R. Consider the
following code and its output.

> mod1 = glm(Ambulance.requests ~ ., data=Ambulance, family=poisson)

> summary(mod1)

Call:
glm(formula = Ambulance.requests ~ ., family = poisson, data = Ambulance)

Deviance Residuals:
Min 1Q Median 3Q Max
-3.2351 -0.8157 -0.0982 0.7787 3.6568

Coefficients:
Estimate Std. Error z value Pr(>|z|)
(Intercept) 2.968477 0.036770 80.732 < 2e-16 ***
WinterYes 0.547756 0.033137 16.530 < 2e-16 ***
WeekendYes -0.607910 0.038184 -15.921 < 2e-16 ***
Bank.holidayYes 0.165684 0.049875 3.322 0.000894 ***
Pollution.levelLow -0.032739 0.042290 -0.774 0.438846
Pollution.levelMedium -0.001587 0.040491 -0.039 0.968734
---
Signif. codes: 0 ‘***’ 0.001 ‘**’ 0.01 ‘*’ 0.05 ‘.’ 0.1 ‘ ’ 1

(Dispersion parameter for poisson family taken to be 1)

Null deviance: 818.08 on 199 degrees of freedom

Residual deviance: 304.97 on 194 degrees of freedom
AIC: 1262.4

Part II, 2018 List of Questions

 105

> mod2 = glm(Ambulance.requests ~ Winter+Weekend, data=Ambulance, family=poisson)

> summary(mod2)

Call:
glm(formula = Ambulance.requests ~ Winter + Weekend, family = poisson,
data = Ambulance)

Deviance Residuals:
Min 1Q Median 3Q Max
-3.4480 -0.8544 -0.1153 0.7689 3.5903

Coefficients:
Estimate Std. Error z value Pr(>|z|)
(Intercept) 2.97077 0.02163 137.34 <2e-16 ***
WinterYes 0.55586 0.03268 17.01 <2e-16 ***
WeekendYes -0.60371 0.03813 -15.84 <2e-16 ***
---
Signif. codes: 0 ‘***’ 0.001 ‘**’ 0.01 ‘*’ 0.05 ‘.’ 0.1 ‘ ’ 1

(Dispersion parameter for poisson family taken to be 1)

Null deviance: 818.08 on 199 degrees of freedom

Residual deviance: 316.39 on 197 degrees of freedom
AIC: 1267.9

Define the two models fitted by this code and perform a hypothesis test with level
1% in which one of the models is the null hypothesis and the other is the alternative. State
the theorem used in this hypothesis test. You may use the information generated by the
following commands.

> qchisq(0.01, df=2, lower.tail=FALSE)

[1] 9.21034
> qchisq(0.01, df=3, lower.tail=FALSE)
[1] 11.34487
> qchisq(0.01, df=4, lower.tail=FALSE)
[1] 13.2767
> qchisq(0.01, df=5, lower.tail=FALSE)
[1] 15.08627

Part II, 2018 List of Questions [TURN OVER

 106

Paper 4, Section II
13J Statistical Modelling
Bridge is a card game played by 2 teams of 2 players each. A bridge club records
the outcomes of many games between teams formed by its m members. The outcomes are
modelled by

exp(βi + βj + β{i,j} − βk − βℓ − β{k,ℓ} )

P(team {i, j} wins against team {k, ℓ}) = ,
1 + exp(βi + βj + β{i,j} − βk − βℓ − β{k,ℓ} )

where βi ∈ R is a parameter representing the skill of player i, and β{i,j} ∈ R is a parameter

representing how well-matched the team formed by i and j is.
(a) Would it make sense to include an intercept in this logistic regression model?
Explain your answer.
(b) Suppose that players 1 and 2 always play together as a team. Is there a unique
maximum likelihood estimate for the parameters β1 , β2 and β{1,2} ? Explain your answer.
(c) Under the model defined above, derive the asymptotic distribution (including the
values of all relevant parameters) for the maximum likelihood estimate of the probability
that team {i, j} wins a game against team {k, ℓ}. You can state it as a function of the true
vector of parameters β, and the Fisher information matrix iN (β) with N games. You may
assume that iN (β)/N → I(β) as N → ∞, and that β has a unique maximum likelihood
estimate for N large enough.

Paper 1, Section II
13J Statistical Modelling
A clinical study follows a number of patients with an illness. Let Yi ∈ [0, ∞) be the
length of time that patient i lives and xi ∈ Rp a vector of predictors, for i ∈ {1, . . . , n}.
We shall assume that Y1 , . . . , Yn are independent. Let fi and Fi be the probability density
function and cumulative distribution function, respectively, of Yi . The hazard function hi
is defined as
fi (t)
hi (t) = for t > 0.
1 − Fi (t)
We shall assume that hi (t) = λ(t) exp(β ⊤ xi ), where β ∈ Rp is a vector of coefficients and
λ(t) is some fixed hazard function.
Rt
(a) Prove that Fi (t) = 1 − exp(− 0 hi (s)ds).
(b) Using the equation in part (a), write the log-likelihood function for β in terms
of λ, β, xi and Yi only.
(c) Show that the maximum likelihood estimate of β can be obtained through a
surrogate Poisson generalised linear model with an offset.

Part II, 2018 List of Questions

 107

Paper 4, Section II
35A Statistical Physics
The one-dimensional Ising model consists of a set of N spins si with Hamiltonian
N
X N
BX
H = −J si si+1 − (si + si+1 ),
2
i=1 i=1

where periodic boundary conditions are imposed so sN +1 = s1 . Here J is a positive

coupling constant and B is an external magnetic field. Define a 2 × 2 matrix M with
elements  
βB
Mst = exp βJst + (s + t) ,
2
where indices s, t take values ±1 and β = (kT )−1 with k Boltzmann’s constant and T
temperature.
(a) Prove that the partition function of the Ising model can be written as

Z = Tr(M N ).

Calculate the eigenvalues of M and hence determine the free energy in the thermodynamic
limit N → ∞. Explain why the Ising model does not exhibit a phase transition in one
dimension.
(b) Consider the case of zero magnetic field B = 0. The correlation function hsi sj i
is defined by
1 X
hsi sj i = si sj e−βH .
Z
{sk }

(i) Show that, for i > 1,

1 X
hs1 si i = st(M i−1 )st (M N −i+1 )ts .
Z s,t

(ii) By diagonalizing M , or otherwise, calculate M p for any positive integer p.

Hence show that

tanhi−1 (βJ) + tanhN −i+1 (βJ)

hs1 si i = .
1 + tanhN (βJ)

Part II, 2018 List of Questions [TURN OVER

 108

Paper 1, Section II
35A Statistical Physics
(a) A macroscopic system has volume V and contains N particles. Let Ω(E, V, N ; δE)
denote the number of states of the system which have energy in the range (E, E + δE)
where δE ≪ E represents experimental uncertainty. Define the entropy S of the system
and explain why the dependence of S on δE is usually negligible. Define the temperature
and pressure of the system and hence obtain the fundamental thermodynamic relation.
(b) A one-dimensional model of rubber consists of a chain of N links, each of length
a. The chain lies along the x-axis with one end fixed at x = 0 and the other at x = L
where L < N a. The chain can “fold back” on itself so x may not increase monotonically
along the chain. Let N→ and N← denote the number of links along which x increases and
decreases, respectively. All links have the same energy.

(i) Show that N→ and N← are uniquely determined by L and N . Determine Ω(L, N ),
the number of different arrangements of the chain, as a function of N→ and N← .
Hence show that, if N→ ≫ 1 and N← ≫ 1 then the entropy of the chain is
    
1 L L
S(L, N ) = kN log 2 − 1+ log 1 +
2 Na Na
   
1 L L
− 1− log 1 −
2 Na Na

where
√ k is Boltzmann’s constant. [You may use Stirling’s approximation: n! ≈
2πn n+1/2 e−n for n ≫ 1.]

(ii) Let f denote the force required to hold the end of the chain fixed at x = L. This
force does work f dL on the chain if the length is increased by dL. Write down
the fundamental thermodynamic relation for this system and hence calculate f as
a function of L and the temperature T .
Assume that N a ≫ L. Show that the chain satisfies Hooke’s law f ∝ L. What
happens if f is held constant and T is increased?

Part II, 2018 List of Questions

 109

Paper 3, Section II
36A Statistical Physics
(a) A system of non-interacting bosons has single particle states |ii with energies
ǫi > 0. Show that the grand canonical partition function is
X  
log Z = − log 1 − e−β(ǫi −µ)
i

where β = 1/(kT ), k is Boltzmann’s constant, and µ is the chemical potential. What is

the maximum possible value for µ?
(b) A system of N ≫ 1 bosons has one energy level with zero energy and M ≫ 1
energy levels with energy ǫ > 0. The number of particles with energies 0, ǫ is N0 , Nǫ
respectively.

(i) Write down expressions for hN0 i and hNǫ i in terms of µ and β.

(ii) At temperature T what is the maximum possible number Nǫmax of bosons

in the state with energy ǫ? What happens for N > Nǫmax ?

(iii) Calculate the temperature TB at which Bose condensation occurs.

(iv) For T > TB , show that µ = ǫ(TB − T )/TB . For T < TB show that

kT eǫ/(kT ) − 1
µ≈− .
N eǫ/(kT ) − eǫ/(kTB )

(v) Calculate the mean energy hEi for T > TB and for T < TB . Hence show
that the heat capacity of the system is
(
1 M ǫ2
kT 2 (eβǫ −1)2 T < TB
C≈ .
0 T > TB

Part II, 2018 List of Questions [TURN OVER

 110

Paper 2, Section II
36A Statistical Physics
(a) Starting from the canonical ensemble, derive the Maxwell–Boltzmann distribu-
tion for the velocities of particles in a classical gas of atoms of mass m. Derive also the
distribution of speeds v of the particles. Calculate the most probable speed.
(b) A certain atom emits photons of frequency ω0 . A gas of these atoms is contained
in a box. A small hole is cut in a wall of the box so that photons can escape in the positive
x-direction where they are received by a detector. The frequency of the photons received
is Doppler shifted according to the formula
 vx 
ω = ω0 1 +
c
where vx is the x-component of the velocity of the atom that emits the photon and c is
the speed of light. Let T be the temperature of the gas.

(i) Calculate the mean value hωi of ω.

p
(ii) Calculate the standard deviation h(ω − hωi)2 i .

(iii) Show that the relative number of photons received with frequency between ω and
ω + dω is I(ω)dω where
I(ω) ∝ exp(−a(ω − ω0 )2 )
for some coefficient a to be determined. Hence explain how observations of the
radiation emitted by the gas can be used to measure its temperature.

Part II, 2018 List of Questions

 111

Paper 4, Section II
29K Stochastic Financial Models
Consider a utility function U : R → R, which is assumed to be concave, strictly
increasing and twice differentiable. Further, U satisfies

U ′ (x) 6 c |x|α , ∀x ∈ R,

for some positive constants c and α. Let X be an N (µ, σ 2 )-distributed random variable
and set f (µ, σ) := E[U (X)].
(a) Show that
   
E U ′ (X)(X − µ) = σ 2 E U ′′ (X) .

∂f ∂f
(b) Show that ∂µ > 0 and ∂σ 6 0. Discuss this result in the context of mean-
variance analysis.
(c) Show that f is concave in µ and σ, i.e. check that the matrix of second derivatives
is negative semi-definite. [You may use without proof the fact that if a 2×2 matrix has non-
positive diagonal entries and a non-negative determinant, then it is negative semi-definite.]

Part II, 2018 List of Questions [TURN OVER

 112

Paper 3, Section II
29K Stochastic Financial Models
Consider a multi-period model with asset prices S̄t = (St0 , . . . , Std ), t ∈ {0, . . . , T },
modelled on a probability space (Ω, F, P) and adapted to a filtration (Ft )t∈{0,...,T } . Assume
that F0 is P-trivial, i.e. P[A] ∈ {0, 1} for all A ∈ F0 , and assume that S 0 is a P-a.s. strictly
positive numéraire, i.e. St0 > 0 P-a.s. for all t ∈ {0, . . . , T }. Further, let Xt = (Xt1 , . . . , Xtd )
denote the discounted price process defined by Xti := Sti /St0 , t ∈ {0, . . . , T }, i ∈ {1, . . . , d}.
(a) What does it mean to say that a self-financing strategy θ̄ is an arbitrage?
(b) State the fundamental theorem of asset pricing.
(c) Let Q be a probability measure on (Ω, F) which is equivalent to P and for which
EQ [|Xt |] < ∞ for all t. Show that the following are equivalent:

(i) Q is a martingale measure.

(ii) If θ̄ = (θ 0 , θ) is self-financing and θ is bounded, i.e. maxt=1,...,T |θt | 6 c < ∞

for a suitable c > 0, then the value process V of θ̄ is a Q-martingale.

(iii) If θ̄ = (θ 0 , θ) is self-financing and θ is bounded, then the value process V

of θ̄ satisfies

EQ [VT ] = V0 .

[Hint: To show that (iii) implies (i) you might find it useful to consider self-financing
strategies θ̄ = (θ 0 , θ) with θ of the form
(
1A if s = t,
θs :=
0 otherwise,

for any A ∈ Ft−1 and any t ∈ {1, . . . , T }.]

(d) Prove that if there exists a martingale measure Q satisfying the conditions in (c)
then there is no arbitrage.

Part II, 2018 List of Questions

 113

Paper 2, Section II
29K Stochastic Financial Models
Consider the Black–Scholes model, i.e. a market model with one risky asset with
price St at time t given by
 
St = S0 exp σBt + µt ,

where (Bt )t>0 denotes a Brownian motion on (Ω, F, P), µ > 0 the constant growth rate,
σ > 0 the constant volatility and S0 > 0 the initial price of the asset. Assume that the
riskless rate of interest is r > 0.
(a) Consider a European option C = f (ST ) with expiry T > 0 for any bounded,
continuous function f : R → R. Use the Cameron–Martin theorem to characterize the
equivalent martingale measure and deduce the following formula for the price πC of C at
time 0:
Z ∞ 
−rT
√
 1 2
πC = e f S0 exp σ T y + (r − 21 σ 2 )T √ e−y /2 dy.
−∞ 2π

(b) Find the price at time 0 of a European option with maturity T > 0 and payoff
C = (ST )γ for some γ > 1. What is the value of the option at any time t ∈ [0, T ]?
Determine a hedging strategy (you only need to specify how many units of the risky asset
are held at any time t).

Part II, 2018 List of Questions [TURN OVER

 114

Paper 1, Section II
30K Stochastic Financial Models
(a) What does it mean to say that (Mn , Fn )n>0 is a martingale?
(b) Let Y1 , Y2 , . . . be independent random variables on (Ω, F, P) with Yi > 0 P-a.s.
and E[Yi ] = 1, i > 1. Further, let
n
Y
M0 = 1 and Mn = Yi , n > 1.
i=1

Show that (Mn )n>0 is a martingale with respect to the filtration Fn = σ(Y1 , . . . , Yn ).
(c) Let X = (Xn )n>0 be an adapted process with respect to a filtration (Fn )n>0
such that E[|Xn |] < ∞ for every n. Show that X admits a unique decomposition

Xn = Mn + An , n > 0,

where M = (Mn )n>0 is a martingale and A = (An )n>0 is a previsible process with A0 = 0,
which can recursively be constructed from X as follows,
 
A0 := 0, An+1 − An := E Xn+1 − Xn | Fn .

(d) Let (Xn )n>0 be a super-martingale. Show that the following are equivalent:

(i) (Xn )n>0 is a martingale.

(ii) E[Xn ] = E[X0 ] for all n ∈ N.

Part II, 2018 List of Questions

 115

Paper 1, Section I
2F Topics in Analysis
State and prove Sperner’s lemma concerning colourings of points in a triangular
grid.
Suppose that △ is a non-degenerate closed triangle with closed edges α1 , α2 and
α3 . Show that we cannot find closed sets Aj with Aj ⊇ αj , for j = 1, 2, 3, such that
3
[ 3
\
Aj = △, but Aj = ∅.
j=1 j=1

Paper 2, Section I
2F Topics in Analysis
For x ∈ Rn we write x = (x1 , x2 , . . . , xn ). Define

P := {x ∈ Rn : xj > 0 for 1 6 j 6 n}.

Qn (a) Suppose that L is a convex Pnsubset of P , that (1, 1, . . . , 1) ∈ L and that

j=1 xj 6 1 for all x ∈ L. Show that j=1 xj 6 n for all x ∈ L.

(b) Suppose that K isQ

a non-empty
Q closed bounded convex subset of P . Show that
there is a u ∈ K such that nj=1 xj 6 nj=1 uj for all x ∈ K. If uj 6= 0 for each j with
1 6 j 6 n, show that
n
X xj
6 n,
uj
j=1

for all x ∈ K, and that u is unique.

Paper 3, Section I
2F Topics in Analysis
State a version of the Baire category theorem and use it to prove the following
result:
If f : C → C is analytic, but not a polynomial, then there exists a point z0 ∈ C
such that each coefficient of the Taylor series of f at z0 is non-zero.

Part II, 2018 List of Questions [TURN OVER

 116

Paper 4, Section I
2F Topics in Analysis
Let 0 6 α < 1 and A > 0. If we have an infinite sequence of integers mn with
1 6 mn 6 Anα , show that
X∞
mn
n!
n=1

is irrational.
Does the result remain true if the mn are not restricted to integer values? Justify
your answer.

Paper 2, Section II
11F Topics in Analysis
(a) Give Bernstein’s probabilistic proof of Weierstrass’s theorem.
(b) Are the following statements true or false? Justify your answer in each case.

(i) If f : R → R is continuous, then there exists a sequence of polynomials Pn

converging pointwise to f on R.

(ii) If f : R → R is continuous, then there exists a sequence of polynomials Pn

converging uniformly to f on R.

(iii) If f : (0, 1] → R is continuous and bounded, then there exists a sequence of

polynomials Pn converging uniformly to f on (0, 1].

(iv) If f : [0, 1] → R is continuous and x1 , x2 , . . . , xm are distinct points in

[0, 1], then there exists a sequence of polynomials Pn with Pn (xj ) = f (xj ),
for j = 1, . . . , m, converging uniformly to f on [0, 1].

(v) If f : [0, 1] → R is m times continuously differentiable, then there exists a

(r)
sequence of polynomials Pn such that Pn → f (r) uniformly on [0, 1] for
each r = 0, . . . , m.

Part II, 2018 List of Questions

 117

Paper 4, Section II
12F Topics in Analysis
We work in C. Consider

K = {z : |z − 2| 6 1} ∪ {z : |z + 2| 6 1}

and
Ω = {z : |z − 2| < 3/2} ∪ {z : |z + 2| < 3/2}.
Show that if f : Ω → C is analytic, then there is a sequence of polynomials pn such that
pn (z) → f (z) uniformly on K.
Show that there is a sequence of polynomials Pn such that Pn (z) → 0 uniformly for
|z − 2| 6 1 and Pn (z) → 1 uniformly for |z + 2| 6 1.
Give two disjoint non-empty bounded closed sets K1 and K2 such that there does
not exist a sequence of polynomials Qn with Qn (z) → 0 uniformly on K1 and Qn (z) → 1
uniformly on K2 . Justify your answer.

Part II, 2018 List of Questions [TURN OVER

 118

Paper 4, Section II
39C Waves
A physical system permits one-dimensional wave propagation in the x-direction
according to the equation
 
∂2 ∂4 ∂2ϕ ∂4ϕ
1−2 2 + 4 + = 0.
∂x ∂x ∂t2 ∂x4

Derive the corresponding dispersion relation and sketch graphs of frequency, phase velocity
and group velocity as functions of the wavenumber. Waves of what wavenumber are at
the front of a dispersing wave train arising from a localised initial disturbance? For waves
of what wavenumbers do wave crests move faster or slower than a packet of waves?
Find the solution of the above equation for the initial disturbance given by
Z ∞
∂ϕ
ϕ(x, 0) = 2A(k)eikx dk , (x, 0) = 0 ,
−∞ ∂t

where A∗ (−k) = A(k), and A∗ is the complex conjugate of A. Let V = x/t be held
fixed. Use the method of stationary phase to√obtain a leading-order approximation to this
solution for large t when 0 < V < Vm = (3 3)/8, where the solutions for the stationary
points should be left in implicit form.
Very briefly discuss the nature of the solutions for −Vm < V < 0 and |V | > Vm .
[Hint: You may quote the result that the large time behaviour of
Z ∞
Φ(x, t) = A(k)eikx−iω(k)t dk ,
−∞

due to a stationary point k = α, is given by

 1
2π 2
Φ(x, t) ∼ ′′
A(α) eiαx−iω(α)t+iσπ/4 ,
|ω (α)| t

where σ = −sgn(ω ′′ (α)). ]

Part II, 2018 List of Questions

 119

Paper 2, Section II
39C Waves
A perfect gas occupies the region x > 0 of a tube that lies parallel to the x-axis.
The gas is initially at rest, with density ρ1 , pressure p1 , speed of sound c1 and specific heat
ratio γ. For times t > 0 a piston, initially at x = 0, is pushed into the gas at a constant
speed V . A shock wave propagates at constant speed U into the undisturbed gas ahead
of the piston. Show that the excess pressure in the gas next to the piston, p2 − p1 ≡ βp1 ,
is given implicitly by the expression

2β 2 p1
V2 = .
2γ + (γ + 1)β ρ1

Show also that

U2 γ+1
=1+ β,
c21 2γ
and interpret this result.
[Hint: You may assume for a perfect gas that the speed of sound is given by
γp
c2 = ,
ρ
and that the internal energy per unit mass is given by
1 p
e= . ]
γ −1ρ

Part II, 2018 List of Questions [TURN OVER

 120

Paper 1, Section II
39C Waves
Derive the wave equation governing the velocity potential for linearised sound waves
in a perfect gas. How is the pressure disturbance related to the velocity potential?
A high pressure gas with unperturbed density ρ0 is contained within a thin metal
spherical shell which makes small amplitude spherically symmetric vibrations. Let the
metal shell have radius a, mass m per unit surface area, and an elastic stiffness which tries
to restore the radius to its equilibrium value a0 with a force κ(a−a0 ) per unit surface area.
Assume that there is a vacuum outside the spherical shell. Show that the frequencies ω of
vibration satisfy  
2 α κa20
θ 1+ = ,
θ cot θ − 1 mc20
where θ = ωa0 /c0 , α = ρ0 a0 /m, and c0 is the speed of sound in the undisturbed gas.
Briefly comment on the existence of solutions.
[Hint: In terms of spherical polar coordinates you may assume that for a function
ψ ≡ ψ(r),
1 ∂2
∇2 ψ = (rψ) . ]
r ∂r 2

Part II, 2018 List of Questions

 121

Paper 3, Section II
40C Waves
Derive the ray-tracing equations
dxi ∂Ω dki ∂Ω dω ∂Ω
= , =− , = ,
dt ∂ki dt ∂xi dt ∂t
for wave propagation through a slowly-varying medium with local dispersion relation
ω = Ω(k; x, t), where ω and k are the frequency and wavevector respectively, t is time and
x = (x, y, z) are spatial coordinates. The meaning of the notation d/dt should be carefully
explained.
A slowly-varying medium has a dispersion relation Ω(k; x, t) = kc(z), where k = |k|.
State and prove Snell’s law relating the angle ψ between a ray and the z-axis to c.
Consider the case of a medium with wavespeed c = c0 (1 + β 2 z 2 ), where β and c0
are positive constants. Show that a ray that passes through the origin with wavevector
k(cos φ, 0, sin φ), remains in the region
 1/2
1 1
|z| 6 zm ≡ −1 .
β | cos φ|

By considering an approximation to the equation for a ray in the region |zm − z| ≪ β −1 ,

or otherwise, determine the path of a ray near zm , and hence sketch rays passing through
the origin for a few sample values of φ in the range 0 < φ < π/2.

Part II, 2018 List of Questions [TURN OVER